A summary of some interesting publications I came accross. Continuously updated.




2023 DeepSpeed ZeRO-Offload++: 6x Higher Training Throughput via Collaborative CPU/GPU Twin-Flow

“System efficiency is still far from optimal when adopting ZeRO-Offload in some scenarios. Especially in the cases like small batch training, model that could not fit into GPU memory but not orders-of-magnitude bigger than GPU memory capacity, CPU offload not only introduce long end-to-end latency, but also underutilize GPU computation resources.” With that in mind, Zero-Offload++ introduces 3 fetures:

  • Twin-Flow: instead having an all-or-nothing policy (ie offload all or none of) in the values to be offloaded, “Twin-Flow allows a portion of optimizer states to be held in CPU memory and the other portion of optimizer states remaining in GPU memory. When optimization step is triggered, both CPU and GPU can do parameter updates simultaneously.” The user can choose the percentage of ratio of parameters in CPU and GPU. “Therefore, with Twin-Flow, we can achieve decent GPU memory and core utilization rate, at the same time reduce training iteation time in optimizer offloading cases.”
  • MemCpy reduction: details not available yet;
  • CPUAdam optimization: details not available yet;


2023 ZeRO++: Extremely Efficient Collective Communication for Giant Model Training, Microsoft

DeepSpeed ZeRO’s compute throughput is limited by the high communication cost from gathering weights in forward pass, backward pass, and averaging gradients. This is mostly prominent on clusters with low-bandwidth, and at very small batch sizes per GPU.

Background, communication pipeline: “Assume the model size as 𝑀. During the forward pass, ZeRO conducts an all-gather operation to collect all the parameters (𝑀) needed to train for all model layers. In the backward pass, ZeRO re-collects parameters (𝑀) with all-gather first, then each GPU can compute local gradients. After that, ZeRO operates reducescatter function to aggregate and redistribute gradients (𝑀) across accelerators. In total, ZeRO has a total communication volume of 3𝑀, spreads evenly across 2 all-gather and 1 reduce-scatter.”

The paper introduces three communication reduction techniques, packed as ZeRO++:

  1. Quantized Weight Communication for ZeRO (qwZ): perform block quantization of the forward all-gather, converting weights from FP16 (2 bytes) to INT8 (1 byte). The main improvement is to replace the typical quantization algorithm (multiplying all parameters by a scalar), by a quantization per block (ie per parameter subset) that includes multiplication by a factor and shifting values by another factor;
  2. Hierarchical Weight Partition for ZeRO (hpZ): data remapping that trades-off communication for more memory and reduces communication overhead of all-gather on weights during backward. Instead of having weights distributed across GPUs, we maintain a full copy on each machine, allowing us to replace the expensive cross-machine all-gather on weights with a faster intra-machine all-gather.
  3. Quantized Gradient Communication for ZeRO (qgZ): replaces the gradients reduce-scatter collective, by doing (1) block-based quantization of gradients to INT4 during communication to reduce the communication size, and recovering the full precision before the reduction operator to preserve training accuracy. Having a fully block-based quantization approach like in (1) was also considered but led to high precision loss and a high error propagation across layers during backpropagation.

The results sections claims that ZeRO++ yields a communication reduction of 4x compared to ZeRO-3, leading to up to 2.16x higher compute throughput on 384 GPUs.


2023 Better speech synthesis through scaling (TorToise), James Bekter

The paper describes a way to apply ML for generating images to the speech synthesis. This result is TorToise, an expressive, multi-voice text-to-speech system. So far, TTS models were hard to train eficiently due to high sampling rate, unavailability of large datasets, or encoder-decoder challenges.

Background: most modern text-to-speech systems operate on speech data that is encoded as a MEL spectrogram. Because of this, most efforts focus on the high-quality decoding of MEL spectrograms back into audio waveforms, a.k.a. a vocoder or a MEL inverter. The author dives in the state-of-art autoregressive transformers and DDPMs models:

  • DALL-E, a transformer model with a (quadratic complexity) full-sequence self-attention, that showed how an autoregressive decoder can be applied to text-toimage generation. The author believes that the “VQVAE decoder used by DALL-E is principally responsible for the blurry incoherence exhibited by most of it’s samples”.
    • DALL-E also introduced the process of re-ranking, that samples from the autoregressive model and picks the best output for downstream use. Re-ranking requires a a strong discriminator to tell good from bad text/image pairings. CLIP was used for this purpose.
  • Denoising diffusion probabilistic models (DDPMs) generate crisp high quality images, and are effective on using low-quality signals to reconstruct the high-dimensional space where those signals derived from. However, DDPMs rely on fixed output shapes, know beforehand. Thus, they “ cannot learn to convert text into audio signals because they cannot solve the implicit alignment problem between text and audio”. Also, DDPMs must be sampled from over multiple iterations, leading to high compute cost and latency.

With that in mind: TorToise works by joining autoregressive decoders and DDPMs: “the autoregressive model will be used to convert a sequence of text tokens to a sequence of tokens representing the output space (in our case, speech tokens). The DDPM will then be used to decode these tokens into a high quality representation of speech.” In practice, for Text-To-Speech, we train the following neural networks:

  • An auto-regressive model on text tokens that yields the probability of each audio token;
  • A contrastive model that ranks outputs of the autoregressive decoder. DALL-E uses CLIP (for images), but TorToise uses Contrastive Language-Voice Pretrained Transformer (CLVP, for TTS).
  • A DDPM to convert speech tokens back into speech spectrograms;

The inputs of the auto-regressive and DDPM models include (or are conditioned to) an additional speech conditioning input, which is one or more audio clips (MEL spectograms) of the same speaker as the target. This allows the model to learn “infer vocal characteristics like tone and prosody” are desired in the target output audio. Finally, they apply the TorToise trick: the DDPM is first trained on converting discrete speech codes into MEL spectrograms, and then fine-tuned on the latent space of the AR model outputs instead of the speech codes. “The logic here is that the AR latent space is far more semantically rich than discrete tokens. By fine-tuning on this latent space, we improve the efficiency of the downstream diffusion model”


2023 Neural Codec Language Models are Zero-Shot Text to Speech Synthesizers (VALL-E), OpenAI

The paper introduces a pipeline for text-to-speech translation (TTS), based on a neural codec language model (VALL-E) using discrete codes (encode/decode embeddings) derived from an off-the-shelf neural audio codec model (Encoded, Défossez et al., 2022). This mode treats TTS as a conditional language modeling task rather than continuous signal regression as in previous work. In practice, contrarily to e.g. AudioLM, a generative audio-to-audio / speech-to-speech model that predicts future audio from input audio, VALL-E is a TTS mode that takes as input a fixed-size text representation and the audio of the first 3 seconds of the text, and tries to predict the future audio that matches the remaining of the input text. VALL-E uses an audio codec code as intermediatte representation and language model as objective, contrary to previous models using mel spectrogram as intermediatte representaion and continuous signal regression as objective. VALL-E is trained with the LibriLight dataset, consisting of 60K hours of English speech with over 7000 unique speakers. This dataset is audio-only, so the authors employ a speech recognition model to generate the (text) transcriptions.

Background, quantization, tokenizer and encoding: audio is typically stored as a sequence of 16-bit integer values, therefore a generative model is required to output \(2^{16}\) = 65536 probabilities per timestep to synthesize the raw audio. Added to the high output size, its long sequence length makes it more intractable for audio synthesis. Therefore, speech quantization is required to compress integer values and sequence length. Common methods are \(\mu\)-law, vector quantization (HuBERT, vq-wav2vec), k-means/self-supervised method, etc. As audio tokenizer, VALL-E uses a pre-trained neural audio codec model, EnCodec, a convolutional encoder-decoder model, whose input and output are both 24 kHz audio across variable bitrates. The encoder produces embeddings at 75 Hz for input waveforms at 24 kHz, which is a 320-fold reduction in the sampling rate. Each embedding is modeled by residual vector quantization (RVQ), with eight hierarchy quantizers with 1024 entries each as shown in Figure 2.

Model architecture: formally speaking, \(Encodec(y) = C^{T \times 8}\), where \(C\) represents the two-dimensional acoustic code matrix (the 8-channel audio embeddings), and \(T\) is the downsample utterance length. Each row in \(C\) represents the eight codes for a given time frame. After quantization, the neural codec decoder is able to reconstruct the waveform, i.e. \(Decodec(C) ≈ \hat{y}\). Given an accoustic prompt matrix \(\hat{C}^{T \times 8}\), the optimization objective of the TTS model is \(max\, p(C \mid x, \hat{C})\), where \(x\) is the corresponding phoneme transcription. I.e. the model learns to extract the content and speaker information from the phoneme sequence and the acoustic prompt, respectively.

There are two models, that refer to the two inference steps:

  1. an auto-regressive (AR) model, a transformer decoder-only architecture, conditioned on the phoneme (text) and accoustic prompt (3-second audio), that gives the discrete tokens of the audio from the first quantizer (Formula 1).
  2. a non auto-regressive (NAR), a transformer decoder will full mask, that regressively predicts the remaining 7 quantizers from the first one (Formula 2).

Note to self: for the use case of synthethising audio in a different language, i.e. that differs from the 3-sec input language and text, see VALL-E X.


2023 High-Fidelity Audio Compression with Improved RVQGAN, Descript Inc.

An audio encoder-decoder that supposedly beats Meta’s encodec. Achieved by combining advances in high-fidelity audio generation with better vector quantization techniques from the image domain, along with improved adversarial and reconstruction losses. Methods:

  • to account for periodicity in audio inputs, they adopted the snake activation function for frequency \(\alpha\) as \(snake(x) = x + \frac{1}{α} sin^2 (αx)\).
  • vanilla VQ-VAEs struggle from low codebook usage due to poor initialization, leading to a significant portion of the codebook being unused. This leads to to poor reconstruction quality. To address this issue, they use two techniques: (1) factorized codes that decouples code lookup and code embedding, by performing code lookup in a low-dimensional space (section 3.2) and (2) L2-normalization of the encoded and codebook vectors converts euclidean distance to cosine similarity, which is helpful for stability and quality.
  • state-of-art applying quantizer dropout degrades the audio reconstruction quality at full bandwidth. To overcome it, they instead apply quantizer dropout to each input example with some probability \(p=0.5\).
  • an improved STFT discriminator at multiple time-scales, that works better in practice and leads to improved phase modeling, compared to Encodec and Soundstream.
  • for frequency domain reconstruction loss, they use a mel-reconstruction loss to improve stability, fidelity and convergence speed; and multi-scale spectral losses to encourage modeling of frequencies in multiple time-scales. For adversarial loss, they use HingeGAN. For codebook learning, they use commitment losses with stop-gradients from the original VQ-VAE formulation. All these losses are weighted to sum up to the final loss.


2023 Llama 2: Open Foundation and Fine-Tuned Chat Model, Meta

LLama 2 is a collection of pretrained and fine-tuned large language models (LLMs) ranging in scale from 7 billion to 70 billion parameters. Llama 2-Chat is a finetuned LLM optimized for dialogue use cases. The models outperform open-source chat models, and based on human evaluations for helpfulness and safety, it outperforms open-source models and appear to be on par with closed-source models (although may not be a suitable substitute). Results on safety human evaluation for Llama 2-Chat are presented in Figure 3. The train dataset is only publicly available sources, which does not include data from Meta’s products or services, or sources that may include users’ personal information. Table 2 presents the GPU compute hours, power consumption and carbon emissions of each model

The pretraining setting and model architecture are adopted from Llama 1, i.e. bytepair encoding (BPE), pre-normalization via RMSNorm, SwiGLU activations, rotary positional embeddings, AdamW optimizer, cosine learning rate scheduler. However, the primary architectural differences from Llama 1 include increased context length and grouped-query attention (GQA).

The finetuning was performed with supervised fine-tuning (Section 3.1), initial and iterative reward modeling (Section 3.2.2) and RLHF (Section 3.2.3). As drawback of RLHF, “initial RLHF models tended to forget the initial instruction after a few turns of dialogue (Figure 9, below, left). To address these limitations, we propose Ghost Attention (GAtt), a very simple method inspired by Context Distillation (Bai et al., 2022b) that hacks the fine-tuning data to help the attention focus in a multi-stage process” (Figure 9, below, right).


2023 LLaMA: Open and Efficient Foundation Language Models, Meta

LLaMa is a collection of Large Language Models (LLM) with 7B to 65B parameters trained in public datasets, with performance superior to GPT-3 and comparable with Chinchilla-70B and PaLM-540B. Training is inspired by the Chinchilla scaling laws. The datasets used for the pre-training data are presented in Table 1, with training hyperparameters in Table 2. String are tokenized using the bytepair encoding (BPE) algorithm, with circa 1.4T tokens after tokenization.

The models architecture is made of several improvements over the original Transformer:

  • Pre-normalization [GPT3]: training stability is improved with RMSNorm normalization at the input of each transformer sub-layer, instead of output.
  • SwiGLU activation function [PaLM]: ReLU activation is replaced with SwiGLU to improve performance, with a dimension of \(\frac{2}{3} 4d\) instead of \(4d\) as in PaLM.
  • Rotary Embeddings [GPTNeo]: positional embeddings are replaced by rotary positional embeddings (RoPE) at each layer of the output.
  • Optimization performed with AdamW optimizer with \(β_1 = 0.9\), \(β2 = 0.95\) and \(eps = 10^{−5}\).
  • Cosine learning rate schedule with a warmup of \(2000\) steps, a weight decay of \(0.1\), a gradient clipping of \(1.0\) and a final learning of \(10%\) of the initial value.
  • Causal multi-Head attention inspired by Rabe and Staats (2021) and uses the backward from Dao et al. (2022), replaces the regular transformer multi-head attention. “This is achieved by not storing the attention weights and not computing the key/query scores that are masked due to the causal nature of the language modeling task.”
  • Activation checkpointing was implemented to reduce memory. Yet it required manually implementing the Pytorch backward propagation function for the Transformer (insted of PyTorch autograd). This also required model and sequence parallelism (why?).
  • Overlap of the computation of activations and the communication between GPUs over the network, to reduce latency.


2023 Sparks of Artificial General Intelligence: Experiments with an early version of GPT-4, Microsoft

A summary paper reporting early results of the experiments with GPT-4 when it was still in active development by OpenAI. The authors “demonstrate that, beyond its mastery of language, GPT-4 can solve novel and difficult tasks that span mathematics, coding, vision, medicine, law, psychology and more, without needing any special prompting. Moreover, in all of these tasks, GPT-4’s performance is strikingly close to human-level performance”. The bulk of the paper contains dozens of examples that compare GPT-4 and Chat-GPT and demonstrate that GPU-4 surpasses ChatGPT in performance, in code generation, audio generation (output as musical notes), drawings (SVG, TIKZ), and mathematical resolutions (LaTeX). As weaknesses, besides the regular hallucinations it was also observed:

  • Incapacity of planning correctly, when planning is not a linear path.
  • Wrong complex arithmetic solver, e.g. What's what is 151412 * 12412 / 12312 + 2? [...] is approximately equal to 152,513.676 instead of 152643.79.
  • Trained on past information only, without temporal awareness, e.g. Whats the president of the US? Donald Trump.
  • lack of rigorous algorithms e.g. What is the 11th letter of "abacadab"? [..] the 11th letter is "b."
  • ilogical reasoning/assumptions due to dataset biases, e.g. in gender: If a man is a computer scientist, a woman is... a source of beauty and inspiration.

But these can be overcome by including external APIs on training and making them in the query e.g.:

  • Using CALC(x) as a calculator of the expression x, what's 151412 * 12412 / 12312 + 2?, or
  • Using SEARCH(x) to search for a query x, who's the president of the US? or
  • Using CALENDAR(subject, date, user) and MAIL(user, text), book a meeting with the title 'subject' on the day 'date' to the user 'user', then email 'user' all the information.


2023 GPTs are GPTs: An Early Look at the Labor Market Impact Potential of Large Language Models, OpenAI, OpenResearch, Univ. of Pennsylvania

The paper investigate the potential implications of large language models (LLMs), on the US labor market. The findings are:

  • around 80% of the U.S. workforce could have at least 10% of their work tasks affected by the introduction of LLMs;
  • approximately 19% of workers may see at least 50% of their tasks impacted;
  • with access to LLMs, about 15% of all worker tasks could be completed significantly faster at the same level of quality; and this share increases to 47%-56% when incorporating software and tooling built on top of LLMs.

The projected effects span all wage levels, with higher-income jobs potentially facing greater exposure to LLM capabilities.


2023 Segment Anything, Meta AI Research

The Segment Anything (SA) project is a task, model, and dataset for image segmentation. SA built the largest segmentation dataset to date (by far), with over 1 billion masks on 11M images. The model is designed to be promptable, to allow zero-shot transter to new image distributions and tasks. The authors clame that its zero-shot performance on new datasets is impressive – competitive or even superior to prior fully supervised results. SA and the dataset are released at https://segment-anything.com


2023 Retentive Network: A Successor to Transformer for Large Language Models, Microsoft and Tsinghua University

(note: a simpler summary video of RetNet can be found here)

RetNet is a multi-scale retention mechanism to substitute multi-head attention in Transformers, which has three computation paradigms:

  • parallel framework, for training parallelism that utilizes GPU devices fully.
  • recurrent framework for low-cost \(O(1)\) inference, which improves decoding throughput (8.4x improvement over Transformer), latency (15.6x), and GPU memory (3.4x) without sacrificing performance, on Figure 1.
  • a chunkwise recurrent representation that can perform efficient long-sequence modeling with linear complexity, where each chunk is encoded parallelly while recurrently summarizing the chunks. It allows encoding each local block for computation speed while recurrently encoding the global blocks to save GPU memory

Retentive network (RetNet) is a stack of \(L\) identical blocks, which follows a similar layout (i.e., residual connection, and pre-LayerNorm) as in Transformer. Each RetNet block contains two modules: a multi-scale retention (MSR) module, and a feed-forward network (FFN) module. The MSR module calls the tokens in a sequence in an auto-regressive manner. The input vector is first created as \(X_0\) in the shape of sequence length by hidden domain size. Then we calculate contextualized vector representations \(X_n\) for each layer of the RetNet. Retention heads can be represented in two alternative ways:

  1. in the parallel representation, where \(Retention(X) = Q K\intercal \dot D)V\) similar to the transformer but with an extra matrix \(D\) (Eq. 5). This is befenicial for parallel training.
  2. in the recurrent representation, it is written as a recurrent neural net (RNN) which is beneficial for inference, and \(Retention(X_n)=Q_n S_n\) where \(S_n\) depends on the previous term \(S_{n-1}\).
  3. a hybrid form combining the previous two representations is also possible to accelerate training on large sequences. Input sequence is divided into chunks. Within each chunk, the computation is performed in the parallel representation. Cross-chunk information is passed in the recurrent representation.

Finally, the model uses \(h = d_{model}/d\) retention heads in each layer, where \(d\) is the head dimension. The heads use different parameter matrices \(W_Q, W_K, W_V \in \mathbb{R}^{d \times d}\) and scalar \(γ\) per head. The overall architecture for a given layer \(l\) of the RetNet is then \(Y_l = MSR(LayerNorm(X_l)) + X_l\) and \(X_{l+1} = FFN(LN(Y_l)) + Y_l\), ie similar to a regular transformer but replacing the attention by a retention head.


2023 Operator Fusion in XLA: Analysis and Evaluation, UToronto

Kernel fusion is the most significant optimization operation in XLA. This paper details XLA and key compiler passes of XLA’s source code. It also presents the speedup that kernel fusion can deliver, and what low-level effects it has on hardware.


2023 LongNet: Scaling Transformers to 1,000,000,000 Tokens, Microsoft and Xi’an Jiaotong University

LongNet is a Transformer variant that can scale the sequence length up to 1B tokens, and without sacrificing the performance on shorter sequences. This overcomes current limitations of attention size in regular transformers, that requires a tradeoff between computational complexity and the model expressivity.

  1. LongNets have a linear computation complexity and a logarithm dependency between any two tokens in a sequence;
  2. They can be served as a distributed trainer for extremely long sequences;
  3. Its main trick is based on the dilated attention, which expands the attentive field exponentially as the distance grows, and is a direct replacement for current attention in Transformers.
    • The general principle behind dilatied attention is: attention allocation decreases exponentially as the distance between tokens grows.
    • When using multiple attention heads, the attention patterns differ among heads by shifting the attention masks.
  4. Dilated attention yields a computation complexity of \(O(ND)\), compared to \(O(N d^2)\) in RNNs, \(O(N^2 d)\) in vanilla attention, and \(O(N \sqrt{N} d)\) in sparse attention.

       


2023 FlashAttention-2: Faster Attention with Better Parallelism and Work Partitioning

“FlashAttention is still not nearly as fast as optimized matrix-multiply (GEMM) operations, reaching only 25-40% of the theoretical maximum FLOPs/s. We observe that the inefficiency is due to suboptimal work partitioning between different thread blocks and warps on the GPU, causing either low-occupancy or unnecessary shared memory reads/writes. We propose FlashAttention-2, with better work partitioning to address these issues. In particular, we (1) tweak the algorithm to reduce the number of non-matmul FLOPs (2) parallelize the attention computation, even for a single head, across different thread blocks to increase occupancy, and (3) within each thread block, distribute the work between warps to reduce communication through shared memory. These yield around 2× speedup compared to FlashAttention, reaching 50-73% of the theoretical maximum FLOPs/s on A100 and getting close to the efficiency of GEMM operations. We empirically validate that when used end-to-end to train GPT-style models, FlashAttention-2 reaches training speed of up to 225 TFLOPs/s per A100 GPU (72% model FLOPs utilization).”


2022 FlashAttention: Fast and Memory-Efficient Exact Attention with IO-Awareness

Transformers are slow and memory-hungry on long sequences, since the time and memory complexity of self-attention are quadratic in sequence length. The authors “ argue that a missing principle is making attention algorithms IOaware—accounting for reads and writes between levels of GPU memory”. To overcome it, Flash Attention improves the attention mechanism with one that uses tiling to reduce the number of memory reads/writes between GPU high bandwidth memory (HBM) and GPU on-chip SRAM. In practice, it tiles the square attention matrix into partial computations that can be computed on the block memory (on-chip SRAM) instead of the global memory (DRAM?). It also trains Transformers faster than existing baselines (15% on bert models and 3x on GPT-2).


2022 TorchScale: Transformers at Scale, Microsoft

TorchScale is an open-source toolkit that allows researchers and developers to scale up Transformers efficiently and effectively. TorchScale adopts Magneto (below) as the default model backbone, thus supporting several applications applications, including language modeling, machine translation, vision pretraining, speech recognition, and multi-modal pretraining, just like Magneto.

  • Stability: to handle models unstability in optimization as the model grows, TorchScale follows the theoretical derivation of DeepNet (below), just like Magneto.
  • Efficiency: TorchScale implements sparsity via Mixture of Experts as in X-MoE , a variant of sparse MoE model. “Torchscale supports both Top-1 and Top-2 routing algorithms (?), which balance the performance and the computation cost. This allows Transformers to scale up to billions or trillions of parameters without much additional computation cost”. Gradient clipping plays an important role in the performance of sparse MoE models, and this was overcome by a new method created by the authors - SparseClip.
    • Note: Gradient clipping is standard practice for deep neural models to alleviate the gradient explosion problem. It is even more important for the sparse MoE models, which are more unstable in training.
  • Transformer variants in the toolkit available at https://github.com/microsoft/torchscale: deep models (DeepNet), Foundation Transformers (Magneto), sparsity ( X-MoE), RetNet, LongNet, parameter stability (SparseClip), …


2022 High Fidelity Neural Audio Compression (Encodec), Meta AI

Encoded is a neural network model for a real-time, high-fidelity, audio codec. It consists in a streaming encoder-decoder architecture with quantized latent space trained in an end-to-end fashion. For faster and simpler training, they use a single multiscale spectrogram that efficiently reduces artifacts and produce high-quality samples. Two main problems arise in lossy neural compression of audio. The first one is overfitting to a subset of audio simples, and it was overcome by using (1) a large and diverse dataset and (2) discriminator networks that serve as perceptual loss. The second problem is compressing efficiently, both in compute time and in size, solved by using residual vector quantization of the neural encoder floating-point output. The authors claim that “designing end-to-end neural compression models is a set of intertwined choices, among which at least the encoder-decoder architecture, the quantization method, and the perceptual loss play key parts”. To that extent, audio quality evaluations (MUSHRA) consist in having humans listen to, compare, and rate excerpts of speech or music compressed with competitive codecs.

Background and model: An audio signal of duration d can be represented by a sequence \(x ∈ [−1, 1]^{C_a × T}\) with \(C_a\) the number of audio channels, \(T = d · f_{sr}\) the number of audio samples at a given sample rate \(f_{sr}\). The EnCodec model is composed of three main components:

  1. an encoder network \(E\) that inputs an audio extract and outputs a latent representation \(z\). It’s simply a stack of 1D convolutions and pooling blocks, followed by a two-layer LSTM for sequence modelling, and a final 1D convolution with \(D\) output channels.
  2. a quantization layer \(Q\) produces a compressed representation \(z_q\), using vector quantization. They use Residual Vector Quantization (RVQ, Zeghidour et al. 2021) to quantize the output of the encoder. As background, general Vector Quantization consists in projecting an input vector onto the closest entry in a codebook of a given size. In this case, RVQ refines this process by computing the residual after quantization, and further quantizing it using a second codebook, and so forth.
  3. a decoder network \(G\) that reconstructs the time-domain signal, \(\hat{x}\), from the compressed latent representation \(z_q\). The decoder’s architecture is the inverse of the encoder, using transposed convolutions instead of strided convolutions.

There are two variants of the model, targetted for the low-latency streamable setup, or a high fidelity non-streamable usage.

The train objective minimizer a linear combination of the following loss products:

  • reconstruction loss, comprised of a time and a frequency domain loss term, to minimize the L1 distance between the target and compressed audio over the time domain, i.e. \(l_t(x, \hat{x}) = \| x − \hat{x} \|_1\). For the frequency domain loss \(l_f\) (Note: type here, this \(l_s\) in the picture above, as in spectogram loss), they use an averaged sum of the L1 and L2 losses between the elements of the input and output mel-spectograms.
  • discriminative loss, a perceptual loss term based on a multi-scale STFT-based (MS-STFT) discriminator, as in Figure 2. In practice, the decoder acts as a generator in a adversarial network and Encodec includes a discriminator module (in orange), with an adversarial loss for the generator as \(l_g(\hat{x}) = \frac{1}{K} \sum_k max (0, 1 − Dk(\hat{x}))\), where \(K\) is the number of discriminators. They add a similar additionally matching loss for the generator, \(l_{feat} (x, \hat{x})\), in Formula 2 (PS: where is this in Figure 1?). The discriminators are trained with the adversarial loss \(l_d(x, \hat{x}) = \frac{1}{K} \sum_{k-1}^K max (0, 1 − D_k(x)) + max(0, 1 + D_k(\hat{x}))\), where \(K\) is the number of discriminators.
  • VQ commitment loss. To support Multi-bandwith learning, at 24 kHz, the model is trained to support the bandwidths 1.5, 3, 6, 12, and 24 kbps by selecting the appropriate number of codebooks to keep in the RVQ step (section 3.2). At 48 kHz, it’s trained to support 3, 6, 12 and 24 kbps. They add a commitment loss \(l_w\) between the output of the encoder, and its quantized value, with no gradient being computed for the quantized value. For each residual step \(c\), where \(C\) is the bandwidth: \(l_w = \sum_{c=1}^C \| z_c - q_c (z_c) \|_2\), where \(q_c(z_c)\) the nearest entry in the corresponding codebook.

The authors also claim that “We introduce a loss balancer in order to stabilize training, in particular the varying scale of the gradients coming from the discriminators” and “We additionally train a small Transformer based language model (Vaswani et al., 2017) with the objective of keeping faster than real time end-to-end compression/decompression on a single CPU core.” (Section 3.3) that I have skipped.


2022 Foundation Transformers (Magneto), Microsoft

Magneto is a “foundation transform for true general-purpose modeling, which serves as a go-to architecture for various tasks and modalities with guaranteed training stability”. Specifically, the paper introduces Sub-LayerNorm (which adds an extra LayerNorm to each sublayer) for good expressivity. Also introduces a novel initialization strategy theoretically derived from DeepNet (below) for stable scaling up. Experiments demonstrate superior performance to the standard Transformer across language, translation, vision, speecha and multimodal tasks. There are only lines of code changes on top of the vanilla Transformer architecture: the Sub-LN structure, and the weights of query projection and key projection are not scaled during initialization following DeepNet (below).


2022 DeepNet: Scaling Transformers to 1,000 Layers

This paper introduces a normalization function (DeepNorm) to modify the residual connection in Transformer, accompanyed with theoretically derived initialization, in order to stabilize extremely deep Transformers.

  • Background: previous work had sown that better initialization methods improve the stability of the training of Transformer.
  • DeepNorm works by introducing a new normalization function at residual connections, which has theoretical justification of bounding the model update by a constant.
  • “The proposed method combines the best of two worlds, i.e., good performance of Post-LN and stable training of Pre-LN (picture above), making DeepNorm a preferred alternative.”.
  • Figure 2 shows the deepnorm (the normalization layer function), deepnorm_init (the weights initialization) and constants.


2022 Contrastive Deep Supervision, Tsinghua University, Intel Corporation, and Xi’an Jiaotong

From the abstract: “the traditional training method only supervises the neural network at its last layer and propagates the supervision layer-by-layer, which leads to hardship in optimizing the intermediate layers. Recently, deep supervision has been proposed to add auxiliary classifiers to the intermediate layers of deep neural networks. By optimizing these auxiliary classifiers with the supervised task loss, the supervision can be applied to the shallow layers directly. However, deep supervision conflicts with the well-known observation that the shallow layers learn low-level features instead of task-biased high-level semantic features. To address this issue, this paper proposes a novel training framework named Contrastive Deep Supervision, which supervises the intermediate layers with augmentation-based contrastive learning”. The rationale is that contrastive learning can provide better supervision for intermediate layers than the supervised task loss. Contrastive learning “regards two augmentations from the same image as a positive pair and different images as negative pairs. During training, the neural network is trained to minimize the distance of a positive pair while maximizing the distance of a negative pair. As a result, the network can learn the invariance to various data augmentation, such as Color Jitter and Random Gray Scale”. Contrastive Deep Supervision starts from those advancements, and optimizes the intermediate layers with contrastive learning instead of traditional supervised learning. As shown in the figure above, “several projection heads are attached in the intermediate layers of the neural networks and trained to perform contrastive learning. These projection heads can be discarded in the inference period to avoid additional computation and storage. Different from deep supervision which trains the intermediate layers to learn the knowledge for a specific task, the intermediate layers in our method are trained to learn the invariance to data augmentation, which makes the neural network generalize better. Besides, since contrastive learning can be performed on unlabeled data, the proposed contrastive deep supervision can also be easily extended in the semi-supervised learning paradigm”. Finally, contrastive deep supervision can be further utilized to boost the performance of knowledge distillation.


2022 Making the Most of Text Semantics to Improve Biomedical Vision-Language Processing, Microsoft

A multi-model approach for text and images applied to health (radioligy), based on contrastive-learning in self-supervised vision-language processing (VLP). As background, one strong motivation is the lack of medical data and the need to self-supervise (and annotate) it, which are expensive and time-consuming. These lead to a interest in multi-model self-supervised learning and cross-model weak supervision, in particular paired text-image data. The paper focus on self-supervised vision-language learning, by jointly learning image and representations for several use cases such as zero-/few-shot image classification, report generation, error detection, and disease localisation.

It introduces a new chest X-ray (CXR) domain-specific language model (CXR-BERT), a self-supervised VLP task for the biomedical use case (BioViL), and a Local Alignment Chest X-ray dataset, MS-CXR. The CXR-BERT is pre-trained with a randomly initialised BERT model via Masked Language Modelling (MLM) (largely following the RoBERTa pretraining configurations), and later fine-tuned with domain-specific data. BioViL uses a convolutional neural network image encoder Eimg, the CXR-BERT text encoder, and projection models to learn representations in a joint space. The CNN model allows them to obtain a grid of local image embeddings, which is fine-grained enough to be useful for segmentation (e.g. 16×16). Each encoder is followed by a modality-specific two-layer perceptron projection model, which projects the encoded modality to a joint space of 128 dimensions. To align the representations and learn a joint embedding, it uses two loss terms based on a symmetric contrastive loss for global alignment of the image and text. After joint training, it uses text prompts to cast the zero-shot classification problem into an image–text similarity task. Results demonstrate CXR-BERT having superior performance and improved vocabulary.


2022 Robust Speech Recognition via Large-Scale Weak Supervision (Whisper), OpenAI

Whisper is an automatic speech recognition (ASR) system, that performs several tasks on a Speech-to-text setup. It is trained on 680K hours of multilingual and multitask supervised data collected from the web. I.e. it’s a weakly supervised dataset. The authors “show that the use of such a large and diverse dataset leads to improved robustness to accents, background noise and technical language.

The Whisper architecture (section 2.2) is an encoder-decoder Transformer. Input audio is split into 30-second chunks, converted into a log-Mel spectrogram, and then passed into an encoder. A decoder is trained to predict the corresponding text caption, intermixed with special tokens that direct the single model to perform tasks such as language identification, phrase-level timestamps, multilingual speech transcription, and to-English speech translation.” Audio is re-sampled to 16,000 Hz, and an 80-channel logmagnitude Mel spectrogram representation is computed on 25-millisecond windows with a stride of 10 milliseconds. For feature normalization, the input is scaled between -1 and 1, and shifted in 0. The encoder and decoder have the same width and number of transformer blocks. Whisper uses the same Byte-Pair Encoding text tokenizer as in GPT2 for english-only models and refit the vocabulary for other languages.

For multi-tasking, they use a simple format to specify all tasks and conditioning information as a sequence of input tokens to the decoder. The beginning of the prediction starts with the <|startoftranscript|> token. In the case where there is no speech in an audio segment, the model is trained to predict a <|nospeech|> token indicating this. The next token specifies the task (either transcription or translation) with an <|transcribe|> or <|translate|>. After this, they specify whether to predict timestamps or not by including a <|notimestamps|> token for that case. At this point, the task and desired format is fully specified, and the output begins. Lastly, they add a <|endoftranscript|> token.

Whisper is evaluated in a zero-shot setting without using any of the training data for each of these datasets so that we are measuring broad generalization. They usei a variant of Word Error Rate to quantify quality. Figure 2 shows improved accuracy over LibriLight-based supervised models on a zero shot setting, and approximate human robustness. I ignored the rest of the paper…


2022 Emergent Abilities of Large Language Models, Google Research & Stanford

The paper discusses instead the phenomenon of emergent abilities of large language models. An ability is emergent if it is not present in smaller models but is present in larger models, and not extrapolated from scaling laws. Phase transition is the scale at which such abilities are exposed. Scale in this context may represent different compute budgets, data quality or other factors - the paper foccuses not on ideal training but on the discussion of such phenomena. As a disclaimer, “model scale is not the singular factor for unlocking an emergent ability” and “the science of training large language models progresses, certain abilities may be unlocked for smaller models with new architectures, higher-quality data, or improved training procedures”.

The first analysis of emergent abilities focuses the prompting paradigm, where outcome is emergent when a model has random performance until a certain scale, after which performance increases to well-above random. This was analysed on 8 different models:

A similar analysis with augmented prompting exposes the emergent property as related to when the model output starts having a positive effect (e.g. being able to do arithmetic only after a certain scale). A multi-step reasoning by providing a chain-of-thoughts as a sequence of intermediatte steps was also analysed, and claimed to be exposed only after \(10^{23}\) FLOPS or approx. 100B parameters. Such scale is also required for intruction following tasks (ie new tasks without prior few-shots exemplars, and only by reading a set of instructions). Program execution tasks require \(9 x 10^{19}\) FLOPS or 40M parameters (for a 8-digit addition) or larger. For model calibration (the ability of a model responding as True of False (or the correctness probability) to which questions they’ll be able to predict correctly) requires \(3*10^{23}\) FLOPS or 52B parameters. It is summarized as:


2022 Rethinking Attention with Performers, Google, Cambridge, DeepMind and Alan Turing Institute

From the abstract: Performers are “Transformer architectures which can estimate regular (softmax) full-rank-attention Transformers with provable accuracy, but using only linear (as opposed to quadratic) space and time complexity, without relying on any priors such as sparsity or low-rankness. To approximate softmax attention kernels, Performers use a novel Fast Attention Via positive Orthogonal Random features approach (FAVOR+)”.

A clearer explanation can be found on this google research post:

Bidirectional attention, where there’s no notion of past and future: by decouplin matrices \(Q′\) and \(K′\) used in lower rank decomposition of \(A\) and conducting matrix multiplications in the order indicated by dashed-boxes, we obtain a linear attention mechanism, never explicitly constructing \(A\) or its approximation:

Left: Standard attention module computation, where the final desired result is computed by performing a matrix multiplication with the attention matrix \(A\) and value tensor \(V\). Right: By decoupling matrices \(Q′\) and \(K′\) used in lower rank decomposition of \(A\) and conducting matrix multiplications in the order indicated by dashed-boxes, we obtain a linear attention mechanism, never explicitly constructing \(A\) or its approximation.

Unidirectional (causal) attention, where tokens do not attend to other tokens appearing later in the sequence: the previous approach is modified to use prefix-sum computations, which only store running totals of matrix computations rather than storing an explicit lower-triangular regular attention matrix.

Left: Standard unidirectional attention requires masking the attention matrix to obtain its lower-triangular part. Right: Unbiased approximation on the LHS can be obtained via a prefix-sum mechanism, where the prefix-sum of the outer-products of random feature maps for keys and value vectors is built on the fly and left-multiplied by query random feature vector to obtain the new row in the resulting matrix.


2022 Training Compute-Optimal Large Language Models, arXiv

Heavily related to HPC’s performance modelling applied to large language models. The authors revisit the question “Given a fixed FLOPs budget, how should one trade-off model size and the number of training tokens?” to which they present three approaches: (1) fix model sizes and vary number of training tokens; (2) vary model sizes for 9 different FLOP counts; (3) fit a parametric loss function to the values retrived from the 2 approaches. Estimates were collected from a total of 400 runs.

The main conclusion is that current large language models are under-performing as they only scaled the model size and not the data size. For compute-optimal training, the model size and number of training tokens should be scalled equally. This hypothesis is demonstrated with a “compute-optimal” model Chinchilla, with the same “compute budget” as Gopher (70B parameters) and 4× more more data. Chinchilla outperforms Gopher (280B), GPT-3 (175B), Jurassic-1 (178B), and Megatron-Turing NLG (530B) on several evaluation tasks.

To be compute optimal (in terms of accuracy vs energy cost), Kaplan et al. (2020) claims that models should not be trained to their lowest possible loss, and for a 10× increase in computational budget, the model should increase by 5.5× and the training tokens by 1.8x. In this paper, the authors defend that model size and training tokens should be scaled in equal proportions.


2021 LoRA: Low-Rank Adaptation of Large Language Models, Microsoft

An important paradigm of natural language processing consists of large-scale pretraining on general domain data and adaptation to particular tasks or domains. Finetuning by retraining all model parameters, can be too expensive or computationally infeasible. Two ways to mitigate this is: (1) to train some (and freeze other) parameters in the model, however falling short of the finetuning baseline; or (2) learning external modules for new tasks, however increasing model latency due to increased depth. Low-Rank Adaptation, or LoRA, freezes the pretrained model weights and injects trainable rank decomposition matrices into each layer of the Transformer architecture, greatly reducing the number of trainable parameters for downstream tasks. “Compared to GPT-3 175B fine-tuned with Adam, LoRA can reduce the number of trainable parameters by 10,000 times and the GPU memory requirement by 3 times”.

Previous work showed that the learned over-parametrized models reside on a low intrinsic dimension. In LoRA, they hypothesize that the change in weights during model adaptation also has a low “intrinsic rank”, i.e. a smaller transformation matrix leading to the proposed Low-Rank Adaptation (LoRA) approach. LoRA trains the dense layers of a neural network using this low-rank matrix, while keeping the layer of the large original model frozen. Moreover, we can perform multi-task adaptation by having a LoRA module per task (ie new A and B matrices below).

It works as follows. Take the gradient matrix \(W_0 \in R^{d × k}\). The gradient update \(W_0 + ∆W\) can be computed instead as \(W_0 + B A\), where \(B ∈ R^{d × r}\), \(A ∈ R^{r×k}\), and the rank \(r \lt \lt min(d, k)\). During training, \(W_0\) is frozen and only \(A\) and \(B\) are updated. The tricky bit: both \(W_0 and ∆W = B A\) are multiplied with the same input, and their respective output vectors are summed coordinate-wise. Therefore, the forward pass (ie the inference operation) is now \(h = W_0x + ∆W x = W_0x + BAx\), so we need to store the unchanged \(W_0\) and trained \(A\) and \(B\). The algorithm is then:

  • \(A\) is initialized with a random gaussian initialization, and \(B\) is initialized as zero, so that \(\Delta W = BA\) at the beginning.
  • We then scale \(∆W x\) by \(\frac{α}{r}\), where \(α\) is a constant in \(r\). The confusing bit: “When optimizing with Adam, tuning α is roughly the same as tuning the learning rate if we scale the initialization appropriately. As a result, we simply set α to the first \(r\) we try and do not tune it.”


2021 Learning Transferable Visual Models From Natural Language Supervision (CLIP), OpenAI

Motivation: State-of-art vision systems are trained on a fixed predetermined set of objects (labels). Additional labeled data is needed to specify any other visual concept. However, “the development of text-to-text as a standardized input-output interface has enabled taskagnostic architectures to zero-shot transfer to downstream datasets, removing the need for specialized output heads or dataset specific customization”. A critical insight is that it is possible to leverage natural language as a flexible prediction space to enable generalization and transfer – ie train a text model, and then specialize it on a non-textual task. In practice, natural language is able to express, and therefore supervise, a much wider set of visual concepts through its generality. Learning from natural language also has an important advantage over most unsupervised or self-supervised learning approaches in that it doesn’t “just” learn a representation but also connects that representation to language which enables flexible zero-shot transfer.

With that in mind, this paper introduces a neural network called CLIP (Contrastive Language–Image Pre-training) which efficiently learns visual concepts from natural language supervision. By design, the network can be instructed in natural language to perform a great variety of classification benchmarks, without directly optimizing for the benchmark’s performance, similar to the “zero-shot” capabilities. CLIP models can then be applied to nearly arbitrary visual classification tasks. Thus, the main keypoint is: by not directly optimizing the model for the benchmark, we show that it becomes much more representative.

In practice, CLIP pre-trains an image encoder and a text encoder to predict which images were paired with which texts in our dataset. The dataset is an abundantly available source of supervision: 400 million pairs of text and respective images found across 500K queries from the internet. For the image encoder, the authors consider 5 differently-sized ResNet50 (with improvements such as attention pooling similar to a QKV attention), and 3 Vision Transformers. The text encoder is a transformer with masked self-attention, with Byte-Pair encoding with a 49152 vocab size. The max sequence length was capped at 76 (section 2.4).

The pipeline for contrastive pre-training (1) is the following:

  • we pass an image through the image encoder (ViT or ResNet). Each image \(i\) is pictured as \(I_i\) above.
  • we pass the text through the text decoder (transformer). Each text \(i\) is pictured as \(T_i\), analogously.
  • the diagonal of the \(I_i T_j\) matrix is the correct image/text labels (in blue).
  • We use contrastive learning and train related text and images to be maximally close, and unrelated to be minimally close. In practice, we maximize the inner product of N labels that go together (full row and full column of the diagonal cells) and minimize the inner product of the N^2-N labels that dont go together (full row and full column of the non-diagonals). We then interpret the product as a logit and then use the softmax on both directions to get the loss (ie it’s a symmetric loss from text and image perspective).
    • in practice, in contrastive learning, “the cosine similarity (i.e. cosine of the angle) of these embeddings is then calculated, scaled by a temperature parameter τ , and normalized into a probability distribution via a softmax”. See figure 3 for source code. scores.
  • In practice, for each image input e.g. \(I_3\), we get the classification distributions \(I_3 T_1\), \(I_3 T_2\),… \(I_3 T_N\).

As you can tell from section (1) in the picture, minibatch size N is crytical. As the minibatch approximates the entire dataset, the representations are more detailed. And the computation of the matrix \(I_i T_j\) increases (quadratically?).

During inference, in (2), we create a dataset by taking a set of labels and adding a prompt to all labels e.g. A photo of a {label}. We then put them through the text encoder, and that is our target set. Then, to perform zero-shot prediciton (3), we take an image, pass it through the image decoder, and gete the classification distribution of that image over the promptedlabels, from where we pick the top label - in the picture A photo of a dog. The main point is: there was zero training needed on the entire task, the image and test datasets can be entirely different, which is a fundamental difference to regular image qualification tasks that have fixed input/output datasets. In CLIP, the model learns the fundamental structure of a language, not just how to difference classes.

Summary of results:

  • Figure 2 shows that “ CLIP is much more efficient at zero-shot transfer than our image caption baseline” and “although highly expressive, they found that transformer-based language models are relatively weak at zero-shot ImageNet classification.”
  • They also do ensembling and combined with prompts, with performs better (Figure 4) by showing better efficiency for the same compute, and equivalent efficient for higher efficiency, is the efficiency/ratio improved.
  • Figure 5 shows that zero-shot CLIP is competitive with fully supervised baselines: in practice, we perform supervising learning of a ResNet model on the ImageNet dataset, then replace the last ResNet layer with a linear layer to allow it to perform a new task. This technique is called Linear Probing method as is based on the fact that the remaining ResNet includes a good representation of the basis. Surprisingly, even on ImageNet where RestNet was trained, the CLIP beats RestNeti50 by +1.9. On the STL10 dataset, the improvement is the highest observed, as this dataset has only a few elements per class, therefore supervised learning is very hard. Similarly, on e.g. MNIST where number of labels is reduced and there are many samples per label, ResNet beats CLIP.
  • In Figure 6, they compare CLIP to few-shot linear probes, and show that CLIP outperforms or matches existing supervised-learning models after linear probing, and that CLIP can improve with linear probing.
  • Following the class count vs accuracy trade-off after linear probing from figure 5, in figure 7 the show the number of labeled examples per class a linear classifier on the same CLIP feature space requires to match the performance of the zero-shot classifier.
  • Figure 9 shows that error goes down as we increase compute and model size. They observed a lot of noise in the results so the conclusions are drawn from the average of all experiments.
  • Figure 10 shows that the CLIP with linear probing beats all state-of-art computer vision models in computer vision tasks, averaged across 12 and 27 datasets.
  • Figure 13 shows the resiliency of the model to perturbation. As in a model is trained on a dataset, and as soon as we change the dataset (but not the labels), the perfomance decreases heavily. The accuracy gap between CLIP and ResNet increases and we degrade the data quality in the dataset (from ImageNet to sketches of ImageNet and adversarial images based on ImageNet).
  • Figure 14, shows that doing linear probe on top of CLIP for e.g. a given dataset, improves accuracy massively for that dataset, but degrades mildly the accuracy of other datasets.
  • Table 7 shows that the prompting matters, by showing that adding the label child to the dataset improves accuracy, dropping the percentage of non-numan or crime-related label assignments dramatically.


2021 ZeRO-Infinity: Breaking the GPU Memory Wall for Extreme Scale Deep Learning, Microsoft

Towards allowing very large models on (memory-limited) GPUs, ZeRO-Infinity is a system technology that leverages GPU, CPU, and NVMe memory to allow for unprecedented model scale on limited resources, without code refactoring. It achieves excellent training throughput and scalability, unencumbered by the limited CPU or NVMe bandwidth. An open source implementation of ZeRO-Infinity is available through DeepSpeed.

ZeRO-Infinity is built on top of ZeRO-3 which partitions all model states to remove memory redundancy, and its main goal is to offload all of the partitioned model states to CPU or NVMe memory, or keep them on the GPU based on the memory requirements. Offloading techniques for different data types are detailed in section 5.

It also introduces a memory-centric tiling for working memory (5.1.3) that breaks a large operator into a sequence of small operators, executed sequentially.


2021 GSPMD: General and Scalable Parallelization for ML Computation Graphs, Google

( also covered on a google blog post )

GSPMD (General and Scalable Parallelization for ML Computation Graphs) is an open-source, automatic, compiler-based parallelization system based on the XLA compiler. Because different model architectures may be better suited to different parallelization strategies, GSPMD is designed to support a large variety of parallelism algorithms appropriate for different use cases (e.g. data parallelism for small models, pipelining parallelism for larger models, or a combination of both).

In GSPMD, each tensor will be assigned a sharding property, either explicitly by the user as initial annotations, or by the sharding completion pass. The sharding property specifies how the data is distributed across devices. GSPMD defines three types of sharding: replicated (all devices have the same full data), tiled (a tiled sharding of the tensor, without data suplication), and partially tilled (an extension to GShard, where tensor is tilled among subgroups of processors, that then have a different tilling within each subgroup).

The sharding properties are user-defined with mesh_split(tensor, device_mesh, dims_mapping) that allows a tensor to be across the device mesh and a mapping from each data tensor dimension (i) to an optional device mesh dimension. This simple API is general enough to express all types of sharding, across the dimension(s) of batch, features, channels and/or others. The automatic partitioner in GSPMD is implemented as transformation/compiler passes in the XLA compiler (Section 3.5), using information about the operator (e.g. dot product is a generalized matrix multiply) or using iterative methods where shardings assigned by the pass are refined incrementally over the iterations.

Left: A simplified feedforward layer of a Transformer model. Blue rectangles represent tensors with dashed red & blue lines overlaid representing the desired partitioning across a 2x2 mesh of devices. Right: A single partition, after GSPMD has been applied. Source: google research post.


2021 Skilful precipitation nowcasting using deep generative models of radar, Google Deepmind

Current weather predictions are done by using numerical weather predictions, by solving physical equations that descrive radar-based wind estimates. Alternative methods use machine learning to capture non-linear behaviour that is not described by the mathematical formalism of the weather-regulating equations. Two main problems arise: poor performance on rarer medium-to-heavy rain events, and weather forecast at high resolution for short-term intervals (2 hours, a.k.a. nowcasting). This paper solves demonstrates improvements in the skill of probabilistic precipitation nowcasting, by using an approach known as generative modelling, based on a deep generative model (DGM) for the probabilistic nowcasting of precipitation.

The DGM presented is a statistical model that learns the probability distributions of data and uses it to generate samples. It also has the ability to simulate many samples from the conditional distribution of future radar given historical radar, thus generating forecasts (similar to ensemble methods). The model predicts \(N\) future radar fields given \(M\) past, using radar-based estimates of surface precipitation \(X_T\) at a time \(T\). Learning is performed similarly to existing work on generative adversarial networks (GANs) for video generation. The GAN model is composed as a generator trained using two discriminators (spatial and temporal) and an additional regularization term. “The generator comprises the conditioning stack which processes past four radar fields that is used as context. […] This stack produces a context representation that is used as an input to the sampler” (that samples from a standard Gaussian).The spatial and temporal discriminator are identical, except that the temporal discrimination uses a 3D kernel to account for the temporal dimension. During evaluation, the generator architecture is the same, but the full radar observations and latent variables distribution of width and heigh \(1/32\) times smaller than then radar observations are used as inputs to the conditioning stack and latent conditioning stack, respectively. “In particular, the latent conditioning stack allows for spatiotemporally consistent predictions for much larger regions than those on which the generator is trained”.

These predictions focus on medium to heavy rain scenarios, as they are the most impactful for society. The model accurately capture large-scale events, while also predicting rainfall uncertainty and generating many alternative rain scenarios (known as ensemble predictions), with consistent predictions of large regions and with lead times from 5–90 min ahead. Results are validated by 50 expert meteorologists that would opt in 89% of situations by this model predictions, compared to competitive methods (PySTEPS, Unet and MetNet) used for nowcasting predictions. As future work, the authors suggest specializing this model for improved long-term predictions.


2021 Reduced, Reused and Recycled: The Life of a Dataset in Machine Learning Research, Google and Univ. California, NeurIPS 2021

Winner of the “Datasets and Benchmarks Best Paper Award” at NeurIPS 2021. Abstract: “We study how dataset usage patterns differ across machine learning subcommunities and across time from 2015-2020. We find increasing concentration on fewer and fewer datasets within task communities, significant adoption of datasets from other tasks, and concentration across the field on datasets that have been introduced by researchers situated within a small number of elite institutions.”


2021 MLP-Mixer: An all-MLP Architecture for Vision, Google, NeurIPS 2021

The paper argues that neither (CNNs) convolution CNNs or (Transformers) attention are necessary for computer vision setups. To that extent, it presents MLP-mixers, a Multi-Layer Perceptron only architecture. “MLP-Mixer contains two types of layers: one with MLPs applied independently to image patches (i.e. “mixing” the per-location features), and one with MLPs applied across patches (i.e. “mixing” spatial information).” Results are competitive with existing methods.


2021 Pay Attention to MLPs, Google, NeurIPS 2021

The paper introduces gMLP (gated MLPs) and show that they can perform as well as Transformers in language and vision applications. It claims that “self-attention is not critical for Vision Transformers, as gMLP can achieve the same accuracy”. In some BERT tasks it performed better than Transformers, and on finetuning tasks, it performed worse (but this can be overcome by making the gMLP model substantially larger).

The gMLPs have no self-attention, and instead rely on channel projections and spatial projections with static parameterization. It consists of a stack of \(L\) blocks with identical size and structure. Each block is defined as:

\[Z = σ(XU), \,\,\,\,\,\,\,\, \tilde{Z} = s(Z), \,\,\,\,\,\,\,\, Y = \tilde{Z} V\]

where \(σ\) is an activation function, \(U\) and \(V\) are linear projections along the channel dimension, and \(s(·)\) is a layer which captures spatial interactions. When \(s\) is an identity mapping, the above transformation degenerates to a regular FFN, ie no cross-token communication. Here, \(s(·)\) is a spatial depthwise convolution (Section 2.1), which, unlike Transformers, does not require position embeddings because that is captured in \(s(·)\).


2021 An Image is Worth 16x16 Words: Transformers for Image Recognition at Scale, Google, ICLR 2021

An extension of the transformer architecture to images. Works by passing as input to the transformer a sequence of linear embeddings of image patches. Paper demonstrates better results on classification tasks, compared to CNNs, ResNets and native attention mechanism (that do not scale well as pixels attend to other pixels leading to a quadratic complexity). Transformers lack the inductive bias of CNNs (e.g. translation equivariance and locality), and therefore do not generalize well when training on insufficient amounts of data. Class is added similarly to BERT as the class token. VTs use 1D positional encodings, since performance of 2D encoders did not deliver significant performance gains. Only MLP layers are local and translationally equivariant, yielding an inductive bias much smaller than CNNs. The hybrid architecture mode uses feature maps of a CNN instead of raw image patches as input. Similar to the original NLP transformer, it scales well and delivers a reduced training time compared to CNN-based architectures. Performance increases with dataset size.


2021 Finetuned Language Models Are Zero-Shot Learners, Google, ICLR 2022

The paper presents a simple method for improving the zero-shot learning abilities of language models. It shows that instruction tuning – finetuning language models on a collection of tasks described via instructions – substantially improves zero-shot performance on unseen tasks. The intuition is that performing instruction tuning—finetuning of the model with datasets expressed via natural language instructions, substantially improves the zero-shot performance of the model. For each dataset, the authors manually compose ten unique templates that use natural language instructions to describe the task for that dataset.


2020 Scaling Laws for Neural Language Models, John Hopkins, OpenAI

Abstract: We study empirical scaling laws for language model performance on the cross-entropy loss. The loss scales as a power-law with model size, dataset size, and the amount of compute used for training, with some trends spanning more than seven orders of magnitude. Other architectural details such as network width or depth have minimal effects within a wide range. Simple equations govern the dependence of overfitting on model/dataset size and the dependence of training speed on model size. These relationships allow us to determine the optimal allocation of a fixed compute budget. Larger models are significantly more sampleefficient, such that optimally compute-efficient training involves training very large models on a relatively modest amount of data and stopping significantly before convergence.

Keypoints:

  • Model performance depends most strongly on scale, which consists of three factors: the number of model parameters N, the size of the dataset D, and the amount of compute C used for training. Performance has a power-law relationship with each of the three scale factors (Fig.1).
  • Within reasonable limits, performance depends very weakly on other architectural hyperparameters such as depth vs. width.
  • Performance improves predictably as long as we scale up N and D in tandem, but enters a regime of diminishing returns if either N or D is held fixed while the other increases.
  • When we evaluate models on text with a different distribution than they were trained on, the results are strongly correlated to those on the training validation set with a roughly constant offset in the loss, i.e. incurs a constant penalty but improves in line with the performance of the training set.
  • When working within a fixed compute budget C but without any other restrictions on the model size N or available data D, we attain optimal performance by training very large models and stopping significantly short of convergence.
  • The ideal batch size for training these models is roughly a power of the loss only


2020 Language Models are Few-Shot Learners (GPT-3), OpenAI

Up until now, substantial gains on many NLP tasks were achieved by pre-training on a large corpus of text followed by fine-tuning on a specific task. This method still requires task-specific fine-tuning datasets of thousands or tens of thousands of examples. This paper shows that scaling up language models greatly improves task-agnostic, few-shot performance, sometimes even reaching competitiveness with prior state-of-the-art finetuning approaches. This paper presents and tests GPT-3 (an autoregressive LLM with 175 billion parameters, 10x larger than previous models) in the few-shot setup.

GPT-3 architecture uses the same model and architecture as GPT-2, including the modified initialization, pre-normalization, and reversible tokenization described therein, with the exception that we use alternating dense and locally banded sparse attention patterns in the layers of the transformer, similar to the Sparse Transformer.

  • Table 2.1 includes the 8 GPT-3 models built and their sizes/hyperparameters.
  • Fig. 2.2 shows the total compute used during training. Based on the analysis in Scaling Laws For Neural Language Models we train much larger models on many fewer tokens than is typical.
  • Fig 3.1 shows the pattern of smooth scaling of performance with compute. Performance (cross-entropy loss) follows a power-law trend with the amount of compute used for training.

Background (Fig 2.1):

  • Fine-Tuning (FT) has been the most common approach in recent years, and involves updating the weights of a pre-trained model by training on a supervised dataset specific to the desired task.
  • Few-Shot (FS) is the term we will use in this work to refer to the setting where the model is given a few demonstrations of the task at inference time as conditioning [RWC+19], but no weight updates are allowed.
  • One-Shot (1S) is the same as few-shot except that only one demonstration is allowed, in addition to a natural language description of the task
  • Zero-Shot (0S) is the same as one-shot except that no demonstrations are allowed, and the model is only given a natural language instruction describing the task.

Tasks tested and performance:

  • On NLP tasks it achieves promising results in the zero-shot and one-shot settings, and in the the few-shot setting is sometimes competitive with or even occasionally surpasses state-of-the-art.
  • It also displays one-shot and few-shot proficiency at tasks designed to test rapid adaption or on-the-fly reasoning, which include unscrambling words, performing arithmetic, and using novel words in a sentence after seeing them defined only once.
  • Fig 3.3 to 3.12 show that GPT3’s performance grows with model size, suggesting that language models continue to absorb knowledge as their capacity increases. Results plotted for the the TriviaQA, translation, Winograd Schema Challenge, PIQA, comprehension, SuperGLUE, ANLI Round 3, arithmetic, word scrambling, and SAT tasks; on zero-, one- and few-shot training, respectively.
  • Fig 3.12 shows that people’s ability to identify whether news articles are model-generated (measured by the ratio of correct assignments to non-neutral assignments) decreases as model size increases. or QuAC.
  • Fig 4.2 plots the benchmark contamination analysis. Data contamination has a minimal effect on GPT-3’s performance on most datasets, but the authors identify a few datasets where it could be inflating results:
  • Chaper 5 details the limitations. GPT-3 struggles with natural language inference tasks like the ANLI dataset, and some reading comprehension datasets like RACE


2020 Graph Transformers Networks, Korea University

One limitation of most GNNs is that they assume the graph structure to be fixed and homogeneous, ie similar types of nodes and edges. From the abstract: “Graph Transformer Networks (GTNs) are capable of generating new graph structures, which involve identifying useful connections between unconnected nodes on the original graph, while learning effective node representation on the new graphs in an end-to-end fashion. Graph Transformer layer, a core layer of GTNs, learns a soft selection of edge types and composite relations for generating useful multi-hop connections”.

  • GTNs perform Meta-Path Generation: a meta-path defines a composite relation \(R = t_1 ◦ t_2 \, ... \, ◦ \, t_l\) between node \(v_1\) and \(v_{l+1}\), where \(R_1 ◦ R_2\) denotes the composition of relation \(R_1\) and \(R_2\).
  • GTNs use graph convolutional network (GCN) to learn useful representations for node classification in an end-to-end fashion.


2019 No Language Left Behind: Scaling Human-Centered Machine Translation, Meta, Berkeley and Johns Hopkins

No Language Left Behing (NLLB) is an open-source projects that provides an ML model capable of delivering igh-quality translations between 200 languages—including low-resource languages like Asturian, Luganda, Urdu and more. The model (NLLB-200) is a conditional model based on Sparsely Gated Mixture of Experts that is “trained on data obtained with novel and effective data mining techniques tailored for low-resource languages”. It also presents “architectural and training improvements to counteract overfitting while training on thousands of tasks”.

  • Background (wikipedia): Mixture of experts (MoE) is a machine learning technique where multiple expert networks (learners) are used to divide a problem space into homogeneous regions. It differs from ensemble techniques in that typically only one or a few expert models will be run, rather than combining results from all models. An example from computer vision is combining one neural network model for human detection with another for pose estimation.


2019 Generating Long Sequences with Sparse Transformers, OpenAI

The paper introduces several sparse factorizations of the attention matrix that reduce the quadratic complexity on memory and runtime Transformers to \(O(n \sqrt{n})\). It also allows for larger sequences. These work by separating the full attention computation into several faster attention operations which, when combined, can approximate the dense attention operation. The authors claim that sparsity in attention is a natural pattern and show (by visual inspection) various examples where most layers had sparse attention patterns across most data points, suggesting that adding sparsity to the attention would not signficantly affecting performance. In other layers, however, they noticed global patterns and data-dependent sparsity, whose performance could be affected by sparsity in the attention matrix.

As a reminder, full self-attention for autoregressive models defines \(S_i = {j : j ≤ i}\), allowing every \(i\)-th element to attend to all previous positions and its own position. In this paper, Factorized self-attention has \(p\) separate attention heads, where the \(m\)-th head defines a subset of the indices \(A_i^{(m)} ⊂ {j : j ≤ i}\) and let \(S_i = A_i^{(m)}\). The problem here is to find efficient choices for the subset \(A\). Section 4.3 details 2D factorization methods via strided attention, or fixed patterns (figure below).

To reduce memory, the authors propose using gradient checkpointing, i.e. recomputing forward-pass attention weights during background pass. To efficiently perform vectorized computation on sparse memory representations, the authors implemented several transpose and slicing operations on the GPU (Section 5.5, note: I’d say they used CUDA shared memory for efficient random memory access).


2019 ALBERT: A Lite BERT for Self-supervised Learning of Language Representations, Google and Toyota

ALBERT (“A Lite BERT”) lowers memory consumption and increase the training speed of BERT. It allows for better scaling, establishes new record performance in several benchmarks, and with fewer parameters than BERT. An ALBERT configuration similar to BERT-large has 18x fewer parameters and can be trained about 1.7x faster. The tecnniques introduced are:

  1. Factorized embedding parameterization, for parameter reduction. “Instead of projecting the one-hot vectors directly into the hidden space of size \(H\), we first project them into a lower dimensional embedding space of size \(E\), and then project it to the hidden space. By using this decomposition, we reduce the embedding parameters from \(O(V × H)\) to \(O(V × E + E × H)\)”. This separation of the size of the hidden layers from the size of vocabulary embedding, makes it easier to grow the hidden size without significantly increasing the parameter size of the vocabulary embeddings.
  2. Cross-layer parameter sharing, for parameter reduction. The authors mention that the parameter reduction also acts as regularisation/generalisation (reduces overfitting as the model learns a representation that generalizes well for all tasks). It does not improve the performance of the model though: “This approach slightly diminishes the accuracy, but the more compact size is well worth the tradeoff”. This technique prevents the parameter from growing with the depth of the network. As a practical example, take a BERT model with 12 layers ie 12 Transformer encoder blocks: instead of learning unique parameters for each layer, ALBERT learns parameters of the first layer reuse the block in the remaining 11 layers.
  3. Self-supervised loss for sentence-order prediction (SOP), for performance improvement. Instead of BERT’s additional loss called next-sentence prediction (NSP, a binary classification loss for predicting whether two segments appear consecutively in the original text), the authors propose SOP, focused on inter-sentence coherence which is designed to address the ineffectiveness of the NSP in BERT. The SOP loss uses as positive examples the same technique as BERT (two consecutive segments from the same document), and as negative examples the same two consecutive segments but with their order swapped. This forces the model to learn finer-grained distinctions about discourse-level coherence properties.


2019 ZeRO: Memory Optimizations Toward Training Trillion Parameter Models, Microsoft

ZeRO (as in Zero Redundancy Optimizer) is a parallelism method that “eliminates memory redundancies in data- and model-parallel training while retaining low communication volume and high computational granularity, allowing us to scale the model size proportional to the number of devices with sustained high efficiency”. The results show the (at the time) largest language model ever created (17B parameters), beating Bert-large (0.3B), GPT-2 (1.5B), Megatron-LM (8.3B), and T5 (11B). It also demonstrates super-linear speedup on 400 GPUs (due to an increase of batch size per accelerator).

As motivation, the authors first emphasize that state-of-art model parallelism splits the model vertically (on each layer), leading to high communication and scaling limitations. Conversely, data parallelism has good compute/communication efficiency but poor memory efficiency. They also analyse “memory consumption of the existing systems on model training and classify it into two parts: 1) For large models, the majority of the memory is occupied by model states which include the optimizer states (such as momentum and variances in Adam), gradients, and parameters. 2) The remaining memory is consumed by activation, temporary buffers and unusable fragmented memory ([…] residual states).” ZeRO-DP claims to have the computation/efficiency of Data Parallelism (DP) while achieving memory efficiency of Model Parallelism (MP). This is achieved by three cumulative optimizations: Optimizer State Partitioning (\(P_{os}\), 4x memory reduction and same communication as DP), Gradient Partitioning (\(P_{os+g}\), 8x memory reduction, same comm.) and Parameter Partitioning (\(P_{os+g+p}\), memory reduction linear with number of accelerations \(N_d\), 50\% increase in communication). ZeRO-DP is at least as memory-efficient and scalable as MP, or more when MP can’t divide the model evenly.

This is achieved by “removing the memory state redundancies across data-parallel processes by partitioning the model states instead of replicating them, and [..] using a dynamic communication schedule during training”. In practice, non-overlapping subsets of layers are delegated to different accelerators. Different optimization levels refer to what content is split or kept across GPUs, as in the figure below. Content that is not replicated but is instead divided in synchronized with dynamic communication across connecting layers. A parameter defining the level of optimization defines the trade-off between variables replicated across accelerators (just like Data Parallelism) and variables split across accelerators (as in Model Parallelism).

At runtime, each processor is allocated a subset of data (DP) and a subset of the model (MP). When that data goes through its layers it will broadcast its layers parameters to other accelerators on the forward pass. Each GPU will run its own data using the received parameters. During the backward pass, gradients will be reduced. See bottom figure and video here.

Finally, ZeRO can be complemented with techniques that reduce activation memory (compression, checkpointing, live analysis). CPU offloading is not recommended or used as “50% of training time can be spent on GPU-CPU-GPU transfers” and this would penalize performance heavily. As a final insight, when compared to MP, “Zero-DP has better scaling efficiency than MP because MP reduces the granularity of the computation while also increasing the communication overhead” and “Zero-R removes the memory redundancies in MP by partitioning the activations checkpoints across GPUs, and uses allgather to reconstruct them on demand”.

ZeRO forward pass: the initial portion of model (\(M_0\)) assigned to \(GPU_0\). It broadcasts its model parameters \(M_0\) to all GPUs (red arrows). Each GPU will do a forward pass of their own data on the received parameters. As we move forward in the model, other GPUs similarly communicate their parameters. The partial activations for each layer are stored by all GPUs. The loss is then computed for each GPU’s data.

image credit: adapted from images in Microsoft Research Blog video

ZeRO backward propagation: on the first iteration of the Backwards pass, GPUs 0,1 and 2 hold the gradients of the last GPU’s model layers \(M_3\) for data points 0, 1 and 2. Combined with the partial activation stored, the partial gradient updates can be computed locally. An all-reduce of all updates will compute the averaged gradient update for model portion \(M_3\) in \(GPU_3\) (green arrows). All remaining layers follow analogously.

image credit: adapted from images in Microsoft Research Blog video


2018 Language Models are Unsupervised Multitask Learners (GPT-2)


2018 Megatron-LM: Training Multi-Billion Parameter Language Models Using Model Parallelism

This paper introduces efficient distributed intra-layer model parallelism for Multi-Layer Perceptrons and Transformer attention mechamism on GPT2, BERT and bidirectional Transformers. The method present is an orthogonal effort and can be combined with existing data, pipeline or model parallelism techniques. The approach approximates closely the pattern of linear scaling (Fig. 1). “We show that the existing BERT architecture results in model degradation as the size increases. We overcome this challenge by rearranging the layer normalization and residual connection in the transformer layers and show that with this change, results for the downstream tasks on development sets improve monotonically as the model size increases.” It also overcomes the limitations of data parallelims (where the model must fit entirely in one worker) and the idleness across time in pipeline parallelism.

Section 3 includes details on the parallelism technique. On the MLP block, take each block being described as \(Y = GeLU(XA)\):

  • the typical approach is to split the weight matrix A along its rows and input X along its columns as (for 2 processors \(1\) and \(2\)): \(X=[X_1, X_2]\) and \(A=[A_1, A_2]^T\). This partitioning will result in \(Y = GeLU(X_1A_1 + X_2A_2)\). Since \(GeLU\) is a nonlinear function, \(GeLU(X_1A_1+ X_2A_2) \neq GeLU(X_1A_1) + GeLU(X_2A_2)\) and this approach will require a synchronization point (to sum both partial sums of products) before the \(GeLU\) function.
  • Another option is to split \(A\) along its columns, i.e. it’s a feature- (not row-) wise partitioning. This allows the \(GeLU\) nonlinearity to be independently applied to the output of each partitioned GEMM: \([Y1, Y2] = [GeLU(XA1), GeLU(XA2)]\). This removes the synchronization point.

This approach splits both GEMMs in the MLP block across GPUs and requires only a single all-reduce operation in the forward pass (g operator) and a single all-reduce in the backward pass (f operator).

The same logic applies to the attention heads, where we split the key, value and query matrices similartly to the matrix \(A\) above. A similar logic follows for the embeddings: “in transformer language models, the output embedding layer shares weights with the input embedding, requiring modifications to both. We parallelize the input embedding weight matrix \(E_{H×v}\) along the vocabulary dimension \(E = [E1, E2]\) (column-wise)”. To reduce the communication in the output of the model (logits) the authors replace communication of logits by scalar losses. Finally, there are 4 total communication operations in the forward and backward pass of a single model parallel transformer layer (Fig. 4).


2018 Averaging Weights Leads to Wider Optima and Better Generalization (Stochastic Weight Averaging), Cornel & Samsumg AI

The authors present SWA, a “simple averaging of multiple points along the trajectory of SGD, with a cyclical or constant learning rate, that leads to better generalization than conventional training” and provides “much flatter solutions than SGD”. The rationale is: (1) SGB with constant or cyclical LR traverse regions of weight space that correspond to high-performing networks, never reaching their central points. (2) Fast Gradient Ensembles (FGE) for \(k\) models required \(l\) times more computation. SWA is an approximation of FGE with the efficiency of a single model, with a better solution that SGD. The algorithm is the following: Starting from \(\hat{w}\) we continue training, using a cyclical or constant learning rate schedule:

  • When using a cyclical learning rate we capture the models \(w_i\) that correspond to the minimum values of the learning rate, i.e. the values at then end of each cycle (at the lowest learning rate value);

  • For high constant learning rates we capture models at each epoch.

Next, we average the weights of all the captured networks wi to get our final model \(w_{SWA}\). For cyclical learning rate schedule, the SWA algorithm is related to FGE, except that instead of averaging the predictions of the models, we average their weights, and we use a different type of learning rate cycle.


2018 GPipe: Efficient Training of Giant Neural Networks using Pipeline Parallelism, Google

GPipe is a method for pipeline parallelism that allows the scaling of neural networks that are expressed as a sequence of layers. The method partitions the original model into subsets of consecutive layers across difference accelerators. This allows for batch pipelining but sequentially feeding to each accelerator one subset of the mini-batch at a time (a micro-batch), and pipelining the whole mini-batch as a sequence of micro-batches. The method applies synchronous (mini-batch) gradient descent with batch accumulation for all micro-batches. During the backward pass, gradients for each micro-batch are computed based on the same model parameters used for the forward pass. At the end of each mini-batch, gradients from all M micro-batches are accumulated and applied to update the model parameters across all accelerators. The authors claim that GPipe’s pipelining (model parallelism) can also be complemented with data parallelism for further training scale. Regular pipelining struggles with RAM issues: when running several micro-batches per mini-batches, it is required to accumulate several input activations (on the forward phase) for the backward phase. Activations (not parameters) are the main memory-consuming factor in CNNs. Therefore, instead of keeping all those activations in memory, “in order to reduce activation memory requirements, […] during forward computation, each accelerator only stores output activations at the partition boundaries, During the backward pass, the accelerator recomputes the composite forward function”. Relating to efficiency and idleness, the “bubble” overhead in the picture can be considered negligible when M ≥ 4 × K, for M micro-batches and K accelerators. “This is also partly because re-computation during the backward pass can be scheduled earlier, without waiting for the gradients from earlier layers”. Benchmark results demonstrate increases peroformance and an almost linear speedup on: image classification (AmoebaNet model) of 480x480 input images, and multilingual translation (128-layer Transformer) tasks. A comparison of runtime against Data Parallelism was not provided. As an important remark, this work was compared with PipeDream that does not follow the Bulk Synchronous Parallel. Moreover, due to the design of overlapping forward and backward passes in PipeDream, it requires maintaining multiple versioned copies of the model parameters. This prevents the PipeDream model to scale as well as GPipe.


2018 PipeDream: Fast and Efficient Pipeline Parallel DNN Training, Microsoft Research, Carnegie Mellon, Stanford

PipeDream is a parallel pipelining method that delivers perfect overlap of communication and computation, and uses all GPUs by overlapping forward and backward passes on data. Compared to other model parallelism techniques, it fully utilises all resources. It “allows perfect overlap of communication and computation. PipeDream keeps all available GPUs productive by systematically partitioning DNN layers among them to balance work and minimize communication, versions model parameters for backward pass correctness, and schedules the forward and backward passes of different inputs in round-robin fashion to optimize time to target accuracy”. On completing the forward pass for a minibatch, each stage asynchronously sends the output activations to the next stage, while simultaneously starting to process another minibatch. Backpropagation proceeds similarly. Thus, the main issue with PipeDream is weight inconsistency (“weight staleness”) caused by performing backward passes of previous mini-batches while doing forward passes of the current mini-batch: “We find that allowing the backward pass for a given minibatch to use more up-to-date parameters than those used in the corresponding forward pass can be a significant problem. PipeDream maintains parameter value versions for each in-flight minibatch to combat this problem”. This leads to an increase of memory requirements. However, it only communicate data between neighboring GPUs, yielding less communication than distributed data parallel, that must communicate all parameters. Finally, PipeDream provides also data parallelism by being able to merge and divide layers across different GPUs. This is supported by: (1) an automatic partitioning scheme to delegate work to compute resources and (2) a work scheduler (“one-forward-one-backward”) that alternates between running a forward and a backward tasks on the queue of tasks available on each GPU to provide a good global flow of the minibatches. A small memory efficiency is achieved by pre-allocating and reusing the GPU memory required for the activations, parameters and intermediate states required in the pipeline, avoiding dynamic allocations. “Experiments with five different DNNs on two different clusters show that PipeDream is up to 5x faster in time-to-accuracy compared to data-parallel training.” Tasks performed: image classification with VGG16 and Inception-v3 models, and video classification with the S2VT model.


2018 BERT: Pre-training of Deep Bidirectional Transformers for Language Understanding, Google

(also covered in a different post )

Existing standard language models are unidirectional and that’s a major limitation in performance, e.g. attending to previous tokens in the self-attention layers in the Transformer. This is an issue for many problems like question answering, it is crucial to incorporate context from both directions. BERT removes this unidirectionality by using a masked language model instead, that allows it to train a deep bidirectional Transformer. BERT model architecture is a multi-layer bidirectional sequence of Transformer encoder blocks. BERT models are trained in 2 steps: pre-training and fine-tuning. During pre-training, the model is trained on unlabeled data on different datasets. During fine-tuned, the pre-trained model is trained for a given specific task. Apart from output layers, the same architectures are used in both pre-training and fine-tuning. During fine-tuning, all parameters are fine-tuned. The input sentence may be a single sentence or a pair of sentences (e.g. question/answer) packed together. Words are embedded with WorkPiece embeddings. [CLS] is the first token of every sentence. [SEP] is a special separator token. To each token (word embedding) it is also added a learned embedding to indicate if it belongs to sentence A or B. Each input is then the sum of its position embedding, segment embedding and token embedding (Fig. 2). The pre-training happens in two unsupervised tasks: (1) Masked LM, by masking of 15% of input tokens at random and trying to predict them, and (2) and Next Sentence Prediction, by passing sentence pairs and predicting whether the second sentence is a logic follow up from the first, or not. The fine-tuning happens differently for every task: we pass the specific inputs and outputs to the BERT and do a regular training. The input is the sequences A and B and separators. The output is the answer to the task by: replacing [CLS] by the sentence or sentence-pair label when the task is to classify a sentence or pair or sentences; replacing the stard and end tokens to indicate the span of output answer tokens that answers the question passed in the input (when input is a question/answer pair, Fig 1); or the class of each word for Named Entity Recognition tasks.


2017 Neural Discrete Representation Learning (RQVAE), Google

The Vector Quantised Variational AutoEncoder (VQ-VAE) aims at learning discrete (not continuous) latent space representations without supervision. It differs from VAEs in two key ways: the encoder network outputs discrete, rather than continuous, codes; and the prior is learnt rather than static. “During forward computation the nearest embedding \(z_q(x)\) (equation 2) is passed to the decoder, and during the backwards pass the gradient \(∇_z\)L is passed unaltered to the encoder. Since the output representation of the encoder and the input to the decoder share the same \(D\) dimensional space, the gradients contain useful information for how the encoder has to change its output to lower the reconstruction loss.” Equation 3 specifies the overall loss function, which has some tricks to allow the model to learn from the discrete mapping of mapping from \(z_e(x)\) to \(z_q(x)\).


2017 Mixed Precision Training, Baidu and NVIDIA

A method for training using half-precision floating point numbers, without losing model accuracy or having to modify hyperparameters. Weights, activations, and gradients are stored in IEEE halfprecision format. This nearly halves memory requirements and speeds up arithmetic. Due to the reduced range of 16- vs 32-bit representation, three techniques are proposed to prevent the loss of critical information (or numerical overflows):

  1. Maintaining a single-precision copy of weights that accumulates the gradients after each optimizer step (this copy is rounded to half-precision for the forward- and back-propagation).
    • Ie, use single-precision for master weights and updates loss-scaling, and accumulating FP16 products into FP32. And half-precision otherwise. This requires memory duplication in two representations for the FP32 master copy of weights: weights, activations and gradients are stored as FP16, and in order to match the accuracy of the FP32 networks, an FP32 master copy of weights is maintained and updated with the weight gradient during the optimizer step.
    • The intuiting behind requiring FP32 for master weights is to include small values that cant be represented by FP16, and because the ratio of the weight value to the weight update is very large.
  2. performing loss-scaling to preserve gradient values with small magnitudes.
    • Rationale: FP16 exponent bias centers the range of normalized value exponents to \([−14, 15]\) while gradient values in practice tend to be dominated by small magnitudes (negative exponents).
    • Scaling up the gradients will shift them to occupy more of the representable range and preserve values that are otherwise lost to zeros
    • To implement scaling, scale the loss value computed in the forward pass by shifting the gradient values into FP16-representable range, prior to back-opagation. By chain rule back-propagation ensures that all the gradient values are scaled by the same amount. Finally, weight gradients must be unscaled before weight update to maintain the update magnitudes as in FP32 training.
    • Any scale factor can be used (without any loss), as long as its product with the maximum absolute gradient value is below 65504.
  3. Using half-precision arithmetic that accumulates into single-precision outputs, which are converted to half precision before storing to memory.
    • Different arithmetics (vector dot-products, reductions, and point-wise operations) require different treatment.
    • “To maintain model accuracy, we found that some networks require that FP16 vector dot-product accumulates the partial products into an FP32 value, which is converted to FP16 before writing to memory. Without this accumulation in FP32, some FP16 models did not match the accuracy of the baseline models.”
    • Large reductions (sums across elements of a vector) e.g. batch-normalization, should be carried out in FP32.
    • Point-wise operations, such as non-linearities and element-wise matrix products, are memory bandwidth limited. Since arithmetic precision does not impact the speed of these operations, either FP16 or FP32 math can be used.


2018 Group Normalization, Facebook AI Research

This paper presents Group normalization. GN surpasses Batch Normalization particularly on small batch sizes, due to error increasing rapidly when the batch size becomes smaller, caused by inaccurate batch statistics estimation. This limits BN’s udage for training larger models and trasferring features. The rationale is that BN exhibits drawbacks that are also caused by its distinct behavior of normalizing along the batch dimension. In particular, it is required for BN to work with a “sufficiently large batch size”.

Layer Normalization and Instance Normalization also avoid normalizing along the batch dimension. These methods are effective for training sequential models (RNN/LSTM) or generative models (GANs), but both have limited success in visual recognition, for which GN presented better results.

Formulation: if \(i = (iN, iC, iH, iW)\) is a 4D vector indexing the image features in (N batch size, Channels , Height, Widht), the mean is \(/mu_i = 1/m \sum_{k \in S_i x_k}\) and the standard deviation is \(\sigma_i = \sqrt{1/m \sum_{k \in S_i} (x_k - \mu_i)^2 }\):

  • Batch norm: \(S_i = \{ k \mid k_C = i_C \}\) ie the output is a vector of the same length as channel count;
  • Layer norm: \(S_i = \{ k \mid k_N = i_N \}\) ie the output is a vector of the same length as batch size;
  • Instance norm: \(S_i = \{ k \mid k_N = i_N, k_C = i_C \}\) ie the output is a matrix of size \(N \times C\);
  • Group norm: \(S_i = \{ k \mid k_N = i_N, \frac{k_C}{C/G} = \frac{i_C}{C/G} \}\) ie the output is a matrix of size \(N \times C/G\), for \(G\) groups;


2016 Semi-Supervised Classification with Graph Convolutional Networks

GCNs are a variant of the Convolutional Neural Networks that operate on graphs. Similarly to CNNs, GCNs learn the features by inspecting neighboring nodes. The main difference is that CNNs are meant to operate on regular Euclidean structures (e.g. images), while GNNs are a generalized applicable to an arbitrary structure or order.

The background is the following. Similarly to the CNN downsampling/upsampling layers, there are several layers in the CGN model. We represent the hidden representation at a given layer as \(H^{(l)}\). We also add the representation of the graph structure in matrix form \(A\), e.g. an adjacency matrix or some function of it. Finally, we write the output of each layer as

\[H^{(l+1)}=f(H^{(l)},A) = σ(AH^{(l)}W^{(l)})\]

with \(H^{(0)}=X\) (i.e. the input data), and where W\(^{(l)}\) is a weight matrix for the \(l\)-th neural network layer and \(σ(⋅)\) is an activation function (e.g. ReLU). To address two limitations (include the node itself on all multiplications with \(A\) and normalizing (by scaling) \(A\), the final formula seen in the paper is then extended to:

\[f(H^{(l)},A)=σ(\hat{D}^{−\frac{1}{2}} \hat{A} \hat{D} ^{−\frac{1}{2}} H^{(l)} W^{(l)}),\]

with \(\hat{A}=A+I\), where \(I\) is the identity matrix and \(\hat{D}\) is the (diagonal) degree matrix of \(\hat{A}\). Note that \(D^{-\frac{1}{2}}\) is the matrix with the reciprocal of the square root of each term in the diagonal. The layer-wise propagation rule is

\[h^{(l+1)}_{v_i}=σ(\sum_j \frac{1}{c_{ij}} h^{(l)}_{v_j} W^{(l)})\]

where \(j\) indexes the neighboring nodes of \(v_i\) and \(c_{ij}\) is a normalization constant for the edge \((v_i,v_j)\).

Sections 2.2 and 2.3 provide theoretical background and section 3 demonstrates an example on the task of node classification, using softmax of the output as in regular CNNs. I used a separate blog post from the author or this post from Francesco Casalegno for a better explanation.


2016 Neural Architecture Search with Reinforcement Learning, Google, ICLR 2017

The authors propose “a recurrent network to generate the model descriptions of neural networks and train this RNN with reinforcement learning to maximize the expected accuracy of the generated architectures on a validation set.” Basically, a DNN that defines the structure of another DNN using RL. The structure and connectivity of the model being designed (the child network) is represented as a variable-length string. This string is generated by the controller network - a recurrent neural network - that uses the child network’s accuracy on the validation set as a reward signal. “The process of generating an architecture stops if the number of layers exceeds a certain value. This value follows a schedule where we increase it as training progresses. Once the controller RNN finishes generating an architecture, a neural network with this architecture is built and trained. At convergence, the accuracy of the network on a held-out validation set is recorded. The parameters of the controller RNN, \(θ_c\), are then optimized in order to maximize the expected validation accuracy of the proposed architectures”. Since the reward signal \(R\) is non-differentiable, they used a policy gradient method to iteratively update \(θ_c\), the REINFORCE rule from Williams (1992). Results beat human-designed networks on the CIFAR-10 dataset and on language modelling with Penn Treebank.


2016 Attention is all you need (Transformer), Google, NeurIPS 2017

(also covered in a different post )

State-of-art transduction models are based on recurrent encoder-decoder architectures (possibly with Attention Mechanisms). The Transformer uses only attention mechanisms, and no recurrence or convolutions. Results show it to be of better performance, more parallelizable (due to non-recurrence in model), and faster to train. Contrarily to recurrent models, the whole source sentence (in the encoder) and target sentence (in the decoder) are fed at once. Therefore, backpropagation happens on a single step as well. Because the concept of word sequence provided by the recurrence was removed, Transformers use positional encoding of the input embeddings based on the combination of sine and cosine waves of different frequencies. The encoder and decoder are composed of a stack of 6 layers each. Each encoder layer includes a multi-heard attention module and a feed forward network. The decoder includes also a third module, a masked multi-head attention, that ensures that sentence does not learn from subsequent words in sentence. An attention head is a mapping of a query to a set of key-value pairs. Key-Value pairs are output by the encoder, and Queries are output by the decoder. The formulation of this dot-product attention is: \(Attention (Q, K, V) = softmax( QK^T / \sqrt{d_k}) V\). Here, the dot-product of all queries and the key (\(QK^T\)) gives a value referring to how well aligned the query vectors are for a given key. This is then converted into a distribution (\(softmax\)) and then used extract the most meaningfull value \(V\) (by multiplying). This is effectively an indexing mechanism (similar to a dictionary \(value = query\_dict[key]\)) but in a continuous space. The scaling factor \(\sqrt{d_k}\) is used to avoid having really small gradients for large values of \(d_k\) (dimensionality of keys). The multi-head attention heads allows the model to jointly attend to information from different (8) representation. It is formulated as \(MultiHead(Q,K, V) = Concat(head_1, ..., head_h)W^O\) where \(head_i = Attention(QW^Q_i ,KW^K_i , VW^V_i)\), ie it’s the linearly-transformed (projected) concatenation of the attention heads with projected Q, K, and V. In terms of performance, self-attention layers have complexity \(O(n^2 d)\) per layer, compared to \(O(n d^2)\) in recurrent models (for sequence length \(n\) and representation dimension \(d\)) - which is typically faster as \(n < d\) in most use cases. It also requires no recurrence and no attention connectivity between previous words in a sentence.


2015 Distilling the Knowledge in a Neural Network, Google, and

2021 Knowledge distillation in deep learning and its applications, and

2020 Knowledge Distillation: A Survey

The distillation method is based on first training a cumbersome model (e.g. an ensemble of models, dropout, etc), and once trained, transfer the knowledge/skill of to a smaller model. In the use case of classification, “an obvious way to transfer the generalization ability of the cumbersome model to a small model is to use the class probabilities produced by the cumbersome model as ‘soft targets’ for training the small model” (while using the same training set). In the use case where the cumbersome model is an ensemble, we can use the geometric mean as target of the small one. The rationale is that “when the soft targets have high entropy, they provide much more information per training case than hard targets and much less variance in the gradients between training cases, so the small model can often be trained on much less data than the original cumbersome model and using a much higher learning rate.” E.g. distillation relies on the fact that the soft max assignments of the large network is a much better label for the input than the Hard MI, thus the smaller network has now the same input trained against a “cleaner” output, and thus requires less complexity to perform equally or better than the larger model .

Applied to the MNIST data case, the previous approach (Caruana et al.) was to use the output of the logits (input to final layer) rather than the probabilities produced by the softmax (as they’re too small), and minimize the squared difference between the logits of the cumbersome and the small model. The authors distillation method is to “raise the temperature of the final softmax until the cumbersome model produces a suitably soft set of targets”, and “then use the same high temperature when training the small model to match these soft targets”.

Few variants of distillation: offline when only the post-training result of the big network is provided to the small network (with the same data, as above), or online when both networks train at the same time. Also possible, for the use case of deep networks, is to use ‘soft labels’ as the output of the bigger network after every X layers, and train the smaller network to learn to replicate the big network’s outputs at every level, and not just the final loss.


2015 Batch Normalization: Accelerating Deep Network Training by Reducing Internal Covariate Shift

quote: Training of DNNs is complicated by the fact that the inputs to each layer are affected by the parameters of all preceding layers – so that small changes to the network parameters amplify as the network becomes deeper. The change in the distributions of layers’ inputs presents a problem because the layers need to continuously adapt to the new distribution. However, […] it has been long known (LeCun et al., 1998b; Wiesler & Ney, 2011) that the network training converges faster if its inputs are whitened – i.e., linearly transformed to have zero means and unit variances, and decorrelated.

Batch Normalization is a technique for training very deep neural networks that standardizes the inputs of the network layers, at every mini-batch. This has the effect of stabilizing the learning process and dramatically reducing the number of training epochs required to train deep networks. Batch Normalization “reduces the internal covariate shift”. “Covariates” is just another name for the input “features”. Covariate shift means the distribution of the features is different in different parts of the training/test data, breaking the i.i.d assumption used across most of ML. This problem occurs frequently in medical data (where you train on samples from one age group, but want to classify something from another age group), or finance (due to changing market conditions). Internal covariate shift refers to covariate shift occurring within a neural network, i.e. going from (say) layer 2 to layer 3. This happens because, as the network learns and the weights are updated, the distribution of outputs of a specific layer in the network changes. This forces the higher layers to adapt to that drift, which slows down learning. BN helps by making the data flowing between intermediate layers of the network look like whitened data, this means you can use a higher learning rate. Since BN has a regularizing effect it also means you can often remove dropout.

In the results, Batch Normalization achieves the same accuracy with 14 times fewer training steps, and beats the original model by a significant margin.


2015 Siamese neural networks for one-shot image recognition, CS Toronto, ICML 2015

The paper describes siamese neural networks for efficient one-shot learning (not to be confused with zero-shot learning), the task of classification under the restriction that we may only observe a single example of each possible class before making a prediction about a test instance. The model learns to identify input pairs according to the probability that they belong to the same class or different classes. This model can then be used to evaluate new images, exactly one per novel class, in a pairwise manner against the test image. The pairing with the highest score according to the verification network is then awarded the highest probability for the one-shot task.

In terms of structure, a siamese neural network consists of twin networks which accept distinct inputs but are joined by an energy function at the top. The parameters between the twin networks are tied (ie the weights are shared), guaranteeing that two extremely similar images could not possibly be mapped by their networks to different locations in feature space because each network computes the same function. I.e. the network is symmetric. For the task of image classification, the authors used a convolutional neural network as the base model of the siamese twins. The units in the final convolutional layer are flattened into a single vector. This convolutional layer is followed by a fully-connected layer, and then one more layer computing the induced distance metric between each siamese twin, which is given to a single sigmoidal output unit. The loss functions is a binary cross-entropy with a regularizer.

The optimization follows a standard backpropagation where the gradient is additive across the twin networks due to the tied weights. Final results show that the model outperforms all available baselines by a significant margin and come close to the best numbers achieved by the previous authors.

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2015 Neural Machine Translation by Jointly Learning to Align and Translate (and Attention Mechanism), D. Bahdanau, K. Cho, Y. Bengio

(also covered in a different post)

In most encoder-decoder models, encoders encode a sentence into a vector of fixed-length, from which a decoder generates the translation. Thus, neural network needs to be able to compress all the necessary information of a source sentence into a fixed-length vector. Here authores claim that fixed-length arrays are a bottleneck in performance on encoder-decoder architectures, particularly for long lentences. Therefore, the authors [quote] “propose to extend this by allowing a model to automatically (soft-)search for parts of a source sentence that are relevant to predicting a target word, without having to form these parts as a hard segment explicitly […] The new architecture consists of a bidirectional RNN as an encoder (BiRNN) and an uni-directional RNN decoder that emulates searching through a source sentence during decoding.”. A BiRNN consists of a forwards a a backward RNNs, containing the summaries of the preceeding words and the following words. The annotation of each word is the concatenation of the forward and backward states. The decoder receives the output of the previous decoded word, a hidden state for time \(i\) (e.g. LSTM hidden state) and the context vector from a sequence of annotations - computed as a weighted sum of annotations. In practice, the encoder encodes the input sentence into a sequence of vectors and the decoder chooses a subset of these vectors adaptively while decoding the translation.


2015 Spatial Transformer Networks, Google DeepMind, NeurIPS 2015


2014 Deeply-supervised Nets, USCD and Microsoft

The rationale of Deeply-supervised nets is the following: in general, a discriminative classifier trained on highly discriminative features will display better performance than a discriminative classifier trained on less discriminative features. If the features in question are the hidden layer feature maps of a deep network, this observation means that the performance of a discriminative classifier trained using these hidden layer feature maps can serve as a proxy for the quality/discriminativeness of those hidden layer feature maps, and further to the quality of the upper layer feature maps. The basic network architecture will be similar to the standard one used in the CNN framework. Our additional deep feedback is brought in by associating a companion local output with each hidden layer. Backpropagation of error now proceeds as usual, with the crucial difference that we now backpropagate not only from the final layer but also simultaneously from our local companion output. Results suggests that it acts as a kind of feature regularization (which leads to significant reduction to the testing error but not necessarily to the train error) and it results in faster convergence, especially in presence of small training data.


2014 Generative Adversarial Networks (GANs), Univ Montreal, NeurIPS 2014

(also detailed on a different blog post)

A new generative model composed of two models trained simultaneously: a generative model G that captures the data distributed, and a discriminative model D that estimates the probability that a sample came from the training data rather than G. The training procedure for G is to maximize the probability of D making a mistake. This framework is the minimax 2-player game. The adversarial framework comes from the generative model facing the discrinative model that learns wether a sample is from the model distribution or the data distribution. “The generative model can be thought of as analogous to a team of counterfeiters, trying to produce fake currency and use it without detection, while the discriminative model is analogous to the police, trying to detect the counterfeit currency. Competition in this game drives both teams to improve their methods until the counterfeits are indistiguishable from the genuine articles.” The generative model generates samples by passing a random noise through a multilayer perceptron. The discriminative model is also a multilayer perceptron. Because both models are connected deep neural networks, training is performed regularly via backpropagation.

image credit: Benjamin Striner, lecture notes CMU 11-785)


2014 Sequence to Sequence Learning with Neural Networks, Google, NeurIPS 2014

(also detailed on a different blog post)

A sequence-to-sequence or Encoder-(to-)Decoder architecture built on Deep Neural Networks of LSTM neurons, demonstrating efficient results on an English-to-French translation task. The main idea is that both Encoder and Decoder are RNNs that use LSTM neurons and its hidden states as a fixed-dimensional vector representation of the sequence so far. That representation is then passed it (i.e. concatenated) to the next token of the sentence. Token [EOS] delimited end of input and output sentences.


2014 Dropout: a simple way to prevent neural networks from overfitting, Univ. Toronto, Journal of ML Research 2014

A method that drops neurons (in different layers) with a given probability \(p\) during train time. For each training minibatch, a new network is sampled. Dropout can be improved by adding max-norm regularization, decaying learning rate and high momentum. At test time, all neurons are used, with outgoing weights multiplied by \(p\). Dropout helps reducing overfitting, as the network learns to never rely on any given activations, so it learns “redundant” ways of solving the task with multiple neurons. It also leads to sparse activations, similar to a regularization (L2). Dropping 20% of input units and 50% of hidden units was often found to be optimal in the original publication. It’s computationally less expensive than regular model averaging of multiple trained DNNs. However, it takes 2-3 times longer to train than single fully-connected DNNs because requires way more epochs, as parameter updates are very noisy. Because a fully connected layer occupies most of the parameters, it is prone to overfitting. Therefore, dropout increases model generalization.


2013 Auto-Encoding Variational Bayes (Variational Autoencoders), Universiteit van Amsterdam, 2013

and An Introduction to Variational Autoencoders from the same authors. Also detailed on a different blog post.

The VAE aims at approximating the distribution of the weights that generates an input, similarly to other variational inference. Te intractable true posterior \(p_{\theta}(z \mid x)\) is approximated by \(q_\phi(z \mid x)\) (the Encoder), whose parameters \(\phi\) are not computed by a closed-form expectation but by the Encoder DNN instead. \(p_\theta(x \mid z)\) is the Decoder, that given a \(z\) will produce/generate the output which is a distribution over the possible values of x. Given a datapoint \(x\) the encoder produces produces a distribution over the possible values of the code \(z\) from which the datapoint \(x\) could have been generated. The VAE proposed includes a DNN decoder, a DNN decoder, with parameters \(\theta\) and \(\phi\), where \(p_\theta(x \mid z)\) is a Gaussian/Bernoulli with distribution parameters computed from \(z\). Therefore the VAE can be viewed as two coupled, independent parameterized models: the encoder/recognition models, and the decoder/generative model (trained together), where the encoder delivers to the decoder an approximation to its posterior over latente random variables. One advantage of the VAE framework, relative to ordinary Variational Inference, is that the encoder is now a (stochastic) function of the input variables, in contrast to VI where each data-case has a separate variational distribution, which is inefficient for large datasets. Finally, the authors noticed that the sampling induces sampling noise in the gradients required for learning (or that because \(z\) is randomly generated and cannot be backpropagated), and to can counteract that variance they use the “reparameterization trick”. It goes as follows: the sample vector \(z\) that is typically sampled from the mean vector \(\mu\) and variance \(\sigma\) in the Gaussian scenario in now described as \(z = \mu + \sigma \cdot \epsilon\) where \(\epsilon\) is always the standard gaussian ie \(\epsilon \sim N(0,1)\).

The loss function is a sum of two terms:

The first term is the reconstruction loss (or expected negative log-likelihood of the i-th datapoint), comparing the model output with the model input and can be the losses we used in the autoencoders(such as L2 loss). The second term is the Kullback-Leibler divergence between the encoder’s distribution \(q_\theta(z\mid x)q\) and \(p(z)\), measuring how much information is lost (in units of nats) when using \(q\) to represent \(p\). It is one measure of how close \(q\) is to \(p\).

VAE vs AE structures. image credit: Data Science Blog: Variational Autoencoders, by Sunil Yadav


2011 Popular Ensemble Methods: An Empirical Study, 2011

A summary of results and conclusions on ensemble methods (bagging, boosting) on DNNs and decision trees. Bagging ensemble generally produces a classifier that is more accurate than a standard classifier. About Boosting: for a few data sets Boosting produced dramatic reductions in error (even compared to Bagging), but for other data sets it actually increases in error over a single classifier (particularly with neural networks). Alternatively, an ensemble of similar DNNs initialized with different random seeds is surprisingly effective, often producing results as good as Bagging. An ideal ensemble consists of highly correct classifiers that disagree as much as possible.

Bagging trains the several different models with different datapoints randomly sampled (with replacement, ie same samples can be redrawn) from the same dataset. Bagging is effective on “unstable” learning algorithms (such as DNNs) where small changes in the training set result in large changes in predictions.

Boosting produces a series of classifiers. The training set used for each member of the series is chosen based on the performance of the earlier classifier(s) in the series. Examples that are incorrectly predicted by previous classifiers in the series are chosen more often than those correctly predicted. Thus Boosting attempts to produce new classifiers that are better able to predict examples for which the current ensemble’s performance is poor. Ada-Boosting can use the approach of (1) selecting a set of examples based on the probabilities of the examples, or (2) simply using all of the examples and weight the error of each example by the probability for that example (i.e., examples with higher probabilities have more effect on the error) – easier as these probabilities are incorporated in the dataset.


2011 Cyclical Learning Rates for Training Neural Networks, US Naval Research Lab, 2017

The author claims that cyclic learning rates improve time to convergence and increases accuracy of most models. It suggests triangular scheduler as a efficient method with similar results to other non-triangular cyclic schedulers. The paper also provides a method to find a good initial learning rate by doing several training short sessions (8 iterations) with different learning rates and picking the best initial learning rate from the analysis. Finally, provides “rule of thumb” parameters for min and max learning rates in the triangular scheduler proposed.


2006 Connectionist Temporal Classification: Labelling Unsegmented: Sequence Data with Recurrent Neural Networks, IDSIA Switzerland, ICML 2006

The paper presents a network and a loss function for the prediction on sequences of labels from unsegmented input data. The overcomes limitations of recursive neural networks that requires well-segmented data. Moreover, it has not been possible to apply RNNs directly to sequence labelling. The problem is that the standard neural network objective functions are defined separately for each point in the training sequence; in other words, RNNs can only be trained to make a series of independent label classifications. The basic idea behind CTC is to interpret the network outputs as a probability distribution over all possible label sequences, conditioned on a given input sequence. Given this distribution, an objective function can be derived that directly maximises the probabilities of the correct labellings. Since the objective function is differentiable, the network can then be trained with standard backpropagation through time. A CTC network has a softmax output layer, with one more unit than there are labels in L. The activations of the first \(L\) units are interpreted as the probabilities of observing the corresponding labels at particular times. The activation of the extra unit is the probability of observing a ‘blank’, or no label. These outputs define the probabilities of all possible ways of aligning all possible label sequences with the input sequence. The total probability of any one label sequence can then be found by summing the probabilities of its different alignments.

The main formulation is of the objective function is: let \(y = N_w(x)\) be the sequence of network outputs, and by \(y^t_k\) the activation of output unit \(k\) at time \(t\), ie the probability of observing label \(k\) at time \(t\). Let \(L'^T\) be the sequence of \(T\) over the \(L' = L ∪ \{blank\}\). \(\, \, \,\) Then: \(p(π \mid x) = \prod_{t=1}^T y^{t}_{π_t}, ∀π ∈ L'^T\). From this equation we observe that that the model assumes frames to be independent. We define the conditional probability of a given labelling \(l ∈ L^{≤T}\) as the sum of the probabilities of all the paths corresponding to it: \(p(l \mid x) = \sum_{\pi} p(\pi \mid x)\). The classified is simply \(h(x) = argmax_l \, p(l \mid x)\).

To efficiently calculate individual labellings, the authors describe the CTC forward-backward algorithm. Training follows the maximum likelihood principle. Experiments compare CTC with framewise method of Hidden Markov Model on the decoding of speech. For fairness, CTC and HMM used the same RNN architecture: bidirectional Long Short-Term Memory. her architecture could have been used instead. BLSTM was chosen because experiments with standard BRNNs and unidirectional networks gave worse results on the same task. Results show improved accuracy of CTC over HMM.


2006 Dimensionality Reduction by Learning an Invariant Mapping (contrastive loss), New York Uni, CVPR 2006

The paper presents Dimensionality Reduction by Learning an Invariant Mapping (DrLIM) - for learning a globally coherent non-linear function that maps the data evenly to the output manifold. The problem is to find a function that maps high dimensional input patterns to lower dimensional outputs, given neighborhood relationships between samples in input space. It also presents the Contrastive Loss Function. The underlying rationale is that a meaningful mapping from high to low dimensional space maps similar input vectors to nearby points on the output manifold and dissimilar vectors to distant points. Therefore, the contrastive loss function runs over pairs of samples. The training is done by (1) collecting images of similar classes using prior knowledge, (2) pair a sample with all other training samples, and (3) traing them against the binary classification (1 or 0) to label them as belonging to the same or different classes, respectively. The neural network used is a “siamese” architecture, consisting of two copies of the function which share the same set of parameters, and a cost module. A loss module whose input is the output of this architecture is placed on top of it. The input to the entire system is a pair of images and a label Y. The images are passed through the functions, yielding two outputs \(G(X_1)\) and \(G(X_2)\). The cost module then generates the distance \(D_W(G_W(X_1), G_W(X_2))\). The loss function combines \(D_W\) with the label to produce the scalar loss \(L_S\) or \(L_D\). The partial loss functions \(L_S\) and \(L_D\) refer to the loss functions to optimize for similar and dissimilar objects. Experiments demonstrate the effectiveness of the method by learning a shift invariant mapping of MNIST samples and a learning temporal neighborhoods and lighting invariance of single objects (airplane).